Electrical World At All

Passive analogue filter development

2009-12-09T22:55:00.000-08:00

Analogue filters are a basic building block of signal processing much used in electronics. Amongst their many applications are the separation of an audio signal before application to bass, mid-range and tweeter loudspeakers; the combining and later separation of multiple telephone conversations onto a single channel; the selection of a chosen radio station in a radio receiver and rejection of others. Passive linear electronic analogue filters are those filters which can be described with linear differential equations (linear); they are composed of capacitors, inductors and, sometimes, resistors (passive) and are designed to operate on continuously varying (analogue) signals. There are many linear filters which are not analogue in implementation (digital filter), and there are many electronic filters which may not have a passive topology – both of which may have the same transfer function of the filters described in this article. Analogue filters are most often used in wave filtering applications, that is, where it is required to pass particular frequency components and to reject others from analogue (continuous-time) signals.
Analogue filters have played an important part in the development of electronics. Especially in the field of telecommunications, filters have been of crucial importance in a number of technological breakthroughs and have been the source of enormous profits for telecommunications companies. It should come as no surprise, therefore, that the early development of filters was intimately connected with transmission lines. Transmission line theory gave rise to filter theory, which initially took a very similar form, and the main application of filters was for use on telecommunication transmission lines. However, the arrival of network synthesis techniques greatly enhanced the degree of control of the designer.Today, it is often preferred to carry out filtering in the digital domain where complex algorithms are much easier to implement, but analogue filters do still find applications, especially for low-order simple filtering tasks and are often still the norm at higher frequencies where digital technology is still impractical, or at least, less cost effective. Wherever possible, and especially at low frequencies, analogue filters are now implemented in a filter topology which is active in order to avoid the wound components required by passive topology.It is possible to design linear analogue mechanical filters using mechanical components which filter mechanical vibrations or acoustic waves. While there are few applications for such devices in mechanics per se, they can be used in electronics with the addition of transducers to convert to and from the electrical domain. Indeed some of the earliest ideas for filters were acoustic resonators because the electronics technology was poorly understood at the time. In principle, the design of such filters can be achieved entirely in terms of the electronic counterparts of mechanical quantities, with kinetic energy, potential energy and heat energy corresponding to the energy in inductors, capacitors and resistors respectively.

Transistor

2009-12-05T08:58:00.000-08:00

A transistor is a semiconductor device commonly used to amplify or switch electronic signals. A transistor is made of a solid piece of a semiconductor material, with at least three terminals for connection to an external circuit. A voltage or current applied to one pair of the transistor's terminals changes the current flowing through another pair of terminals. Because the controlled (output) power can be much more than the controlling (input) power, the transistor provides amplification of a signal. Some transistors are packaged individually but most are found in integrated circuits.
History:-
Physicist Julius Edgar Lilienfeld filed the first patent for a transistor in Canada in 1925, describing a device similar to a Field Effect Transistor or "FET".However, Lilienfeld did not publish any research articles about his devices,[citation needed] and in 1934, German inventor Oskar Heil patented a similar device.In 1947, John Bardeen and Walter Brattain at AT&T's Bell Labs in the United States observed that when electrical contacts were applied to a crystal of germanium, the output power was larger than the input. Solid State Physics Group leader William Shockley saw the potential in this, and over the next few months worked to greatly expand the knowledge of semiconductors, and thus could be described as the "father of the transistor". The term was coined by John R. Pierce. According to physicist/historian Robert Arns, legal papers from the Bell Labs patent show that William Shockley and Gerald Pearson had built operational versions from Lilienfeld's patents, yet they never referenced this work in any of their later research papers or historical articles.The first silicon transistor was produced by Texas Instruments in 1954.[5] This was the work of Gordon Teal, an expert in growing crystals of high purity, who had previously worked at Bell Labs. The first MOS transistor actually built was by Kahng and Atalla at Bell Labs in 1960.
Importants:-
The transistor is considered by many to be one of the greatest inventions of the twentieth century. The transistor is the key active component in practically all modern electronics. Its importance in today's society rests on its ability to be mass produced using a highly automated process (fabrication) that achieves astonishingly low per-transistor costs.Although several companies each produce over a billion individually-packaged (known as discrete) transistors every year, the vast majority of transistors produced are in integrated circuits (often shortened to IC, microchips or simply chips) along with diodes, resistors, capacitors and other electronic components to produce complete electronic circuits. A logic gate consists of up to about twenty transistors whereas an advanced microprocessor, as of 2006, can use as many as 1.7 billion transistors (MOSFETs). "About 60 million transistors were built this year [2002] ... for [each] man, woman, and child on Earth.The transistor's low cost, flexibility, and reliability have made it a ubiquitous device. Transistorized mechatronic circuits have replaced electromechanical devices in controlling appliances and machinery. It is often easier and cheaper to use a standard microcontroller and write a computer program to carry out a control function than to design an equivalent mechanical control function.
Uses:-
The bipolar junction transistor, or BJT, was the most commonly used transistor in the 1960s and 70s. Even after MOSFETs became widely available, the BJT remained the transistor of choice for many analog circuits such as simple amplifiers because of their greater linearity and ease of manufacture. Desirable properties of MOSFETs, such as their utility in low-power devices, usually in the CMOS configuration, allowed them to capture nearly all market share for digital circuits; more recently MOSFETs have captured most analog and power applications as well, including modern clocked analog circuits, voltage regulators, amplifiers, power transmitters, motor drivers, etc.

Electrical impedance

2009-12-05T08:55:00.000-08:00

Electrical impedance, or simply impedance, describes a measure of opposition to alternating current (AC). Electrical impedance extends the concept of resistance to AC circuits, describing not only the relative amplitudes of the voltage and current, but also the relative phases. When the circuit is driven with direct current (DC) there is no distinction between impedance and resistance; the latter can be thought of as impedance with zero phase angle.The symbol for impedance is usually \scriptstyle Z and it may be represented by writing its magnitude and phase in the form \scriptstyle Z \angle \theta . However, complex number representation is more powerful for circuit analysis purposes. The term impedance was coined by Oliver Heaviside in July 1886.Arthur Kennelly was the first to represent impedance with complex numbers in 1893.Impedance is defined as the frequency domain ratio of the voltage to the current. In other words, it is voltage–current ratio for a single complex exponential at a particular frequency ω. In general, impedance will be a complex number, but this complex number has the same units as resistance, for which the SI unit is the ohm. For a sinusoidal current or voltage input, the polar form of the complex impedance relates the amplitude and phase of the voltage and current. In particular,

Hall effect sensor

2009-12-04T08:18:00.000-08:00

A Hall effect sensor is a transducer that varies its output voltage in response to changes in magnetic field. Hall sensors are used for proximity switching, positioning, speed detection, and current sensing applications.In its simplest form, the sensor operates as an analogue transducer, directly returning a voltage. With a known magnetic field, its distance from the Hall plate can be determined. Using groups of sensors, the relative position of the magnet can be deduced.Electricity carried through a conductor will produce a magnetic field that varies with current, and a Hall sensor can be used to measure the current without interrupting the circuit. Typically, the sensor is integrated with a wound core or permanent magnet that surrounds the conductor to be measured.Frequently, a Hall sensor is combined with circuitry that allows the device to act in a digital (on/off) mode, and may be called a switch in this configuration. Commonly seen in industrial applications such as the pictured pneumatic cylinder, they are also used in consumer equipment; for example some computer printers use them to detect missing paper and open covers. When high reliability is required, they are used in keyboards.Hall sensors are commonly used to time the speed of wheels and shafts, such as for internal combustion engine ignition timing or tachometers. They are used in brushless DC electric motors to detect the position of the permanent magnet. In the pictured wheel carrying two equally spaced magnets, the voltage from the sensor will peak twice for each revolution. This arrangement is commonly used to regulate the speed of disc drives.
History:-
A hall probe contains an indium compound crystal mounted on an aluminum backing plate, and encapsulated in the probe head. The plane of the crystal is perpendicular to the probe handle. Connecting leads from the crystal are brought down through the handle to the circuit box.When the Hall Probe is held so that the magnetic field lines are passing at right angles through the sensor of the probe, the meter gives a reading of the value of magnetic flux density (B). A current is passed through the crystal which, when placed in a magnetic field has a “Hall Effect” voltage developed across it. The Hall Effect is seen when a conductor is passed through a uniform magnetic field. The natural electron drift of the charge carriers causes the magnetic field to apply a Lorentz force (the force exerted on a charged particle in an electromagnetic field) to these charge carriers. The result is what is seen as a charge separation, with a build up of either positive or negative charges on the bottom or on the top of the plate. The crystal measures 5 mm square. The probe handle, being made of a non-ferrous material, has no disturbing effect on the field.A Hall Probe is sensitive enough to measure the Earth's magnetic field. It must be held so that the Earth's field lines are passing directly through it. It is then rotated quickly so the field lines pass through the sensor in the opposite direction. The change in the flux density reading is double the Earth's magnetic flux density. A hall probe must first be calibrated against a known value of magnetic field strength. For a solenoid the hall probe is placed in the center

PID controller

2009-12-03T08:09:00.000-08:00

A proportional–integral–derivative controller (PID controller) is a generic control loop feedback mechanism (controller) widely used in industrial control systems. A PID controller attempts to correct the error between a measured process variable and a desired setpoint by calculating and then instigating a corrective action that can adjust the process accordingly and rapidly, to keep the error minimal.
General:-
The PID controller calculation (algorithm) involves three separate parameters; the proportional, the integral and derivative values. The proportional value determines the reaction to the current error, the integral value determines the reaction based on the sum of recent errors, and the derivative value determines the reaction based on the rate at which the error has been changing. The weighted sum of these three actions is used to adjust the process via a control element such as the position of a control valve or the power supply of a heating element.By tuning the three constants in the PID controller algorithm, the controller can provide control action designed for specific process requirements. The response of the controller can be described in terms of the responsiveness of the controller to an error, the degree to which the controller overshoots the setpoint and the degree of system oscillation. Note that the use of the PID algorithm for control does not guarantee optimal control of the system or system stability.
Some applications may require using only one or two modes to provide the appropriate system control. This is achieved by setting the gain of undesired control outputs to zero. A PID controller will be called a PI, PD, P or I controller in the absence of the respective control actions. PI controllers are particularly common, since derivative action is very sensitive to measurement noise, and the absence of an integral value may prevent the system from reaching its target value due to the control action.
Note: Due to the diversity of the field of control theory and application, many naming conventions for the relevant variables are in common use.
History:-
A familiar example of a control loop is the action taken to keep one's shower water at the ideal temperature, which typically involves the mixing of two process streams, cold and hot water. The person feels the water to estimate its temperature. Based on this measurement they perform a control action: use the cold water tap to adjust the process. The person would repeat this input-output control loop, adjusting the hot water flow until the process temperature stabilized at the desired value.Feeling the water temperature is taking a measurement of the process value or process variable (PV). The desired temperature is called the setpoint (SP). The output from the controller and input to the process (the tap position) is called the manipulated variable (MV). The difference between the measurement and the setpoint is the error (e), too hot or too cold and by how much.As a controller, one decides roughly how much to change the tap position (MV) after one determines the temperature (PV), and therefore the error. This first estimate is the equivalent of the proportional action of a PID controller. The integral action of a PID controller can be thought of as gradually adjusting the temperature when it is almost right. Derivative action can be thought of as noticing the water temperature is getting hotter or colder, and how fast, anticipating further change and tempering adjustments for a soft landing at the desired temperature (SP).Making a change that is too large when the error is small is equivalent to a high gain controller and will lead to overshoot. If the controller were to repeatedly make changes that were too large and repeatedly overshoot the target, the output would oscillate around the setpoint in either a constant, growing, or decaying sinusoid. If the oscillations increase with time then the system is unstable, whereas if they decay the system is stable. If the oscillations remain at a constant magnitude the system is marginally stable. A human would not do this because we are adaptive controllers, learning from the process history, but PID controllers do not have the ability to learn and must be set up correctly. Selecting the correct gains for effective control is known as tuning the controller.If a controller starts from a stable state at zero error (PV = SP), then further changes by the controller will be in response to changes in other measured or unmeasured inputs to the process that impact on the process, and hence on the PV. Variables that impact on the process other than the MV are known as disturbances. Generally controllers are used to reject disturbances and/or implement setpoint changes. Changes in feed water temperature constitute a disturbance to the shower process.In theory, a controller can be used to control any process which has a measurable output (PV), a known ideal value for that output (SP) and an input to the process (MV) that will affect the relevant PV. Controllers are used in industry to regulate temperature, pressure, flow rate, chemical composition, speed and practically every other variable for which a measurement exists. Automobile cruise control is an example of a process which utilizes automated control.
Due to their long history, simplicity, well grounded theory and simple setup and maintenance requirements, PID controllers are the controllers of choice for many of these applications.

Inductance

2009-12-02T00:02:00.000-08:00

Inductance is the property in an electrical circuit where a change in the electric current through that circuit induces an electromotive force (EMF) that opposes the change in current (See Induced EMF).In electrical circuits, any electric current i produces a magnetic field and hence generates a total magnetic flux Φ acting on the circuit. This magnetic flux, due to Lenz's law, tends to act to oppose changes in the flux by generating a voltage (a back EMF) in the circuit that counters or tends to reduce the rate of change in the current. The ratio of the magnetic flux to the current is called the self-inductance, which is usually simply referred to as the inductance of the circuit. To add inductance to a circuit, electronic components called inductors are used, which consist of coils of wire to concentrate the magnetic field.The term 'inductance' was coined by Oliver Heaviside in February 1886.It is customary to use the symbol L for inductance, possibly in honour of the physicist Heinrich Lenz.

Electromagnetic spectrum

2009-11-27T08:47:00.000-08:00

The electromagnetic spectrum is the range of all possible frequencies of electromagnetic radiation. The "electromagnetic spectrum" of an object is the characteristic distribution of electromagnetic radiation emitted or absorbed by that particular object.The electromagnetic spectrum extends from below frequencies used for modern radio to gamma radiation at the short-wavelength end, covering wavelengths from thousands of kilometers down to a fraction of the size of an atom. The long wavelength limit is the size of the universe itself, while it is thought that the short wavelength limit is in the vicinity of the Planck length, although in principle the spectrum is infinite and continuous.
Type:-
While the classification scheme is generally accurate, in reality there is often some overlap between neighboring types of electromagnetic energy. For example, radio waves at 60 Hz may be received and studied by astronomers, or may be ducted along wires as electric power.
The distinction between X and gamma rays is based on sources: gamma rays are the photons generated from nuclear decay or other nuclear and /particle process, whereas X-rays are generated by electronic transitions involving highly energetic inner atomic electrons.
Also, the region of the spectrum of the particular electromagnetic radiation is reference-frame dependent (on account of the Doppler shift for light) so EM radiation which one observer would say is in one region of the spectrum could appear to an observer moving at a substantial fraction of the speed of light with respect to the first to be in another part of the spectrum. For example, consider the cosmic microwave background. It was produced, when matter and radiation decoupled, by the de-excitation of hydrogen atoms to the ground state. These photons were from Lyman series transitions, putting them in the ultraviolet (UV) part of the electromagnetic spectrum. Now this radiation has undergone enough cosmological red shift to put it into the microwave region of the spectrum for observers moving slowly (compared to the speed of light) with respect to the cosmos. However, for particles moving near the speed of light, this radiation will be blue-shifted in their rest frame. The highest energy cosmic ray protons are moving such that, in their rest frame, this radiation is to high energy gamma rays which interact with the proton to produce bound pairs . This is the source of the limit.

Direct current

2009-11-22T04:35:00.000-08:00

Direct current (DC) is the unidirectional flow of electric charge. Direct current is produced by such sources as batteries, thermocouples, solar cells, and commutator-type electric machines of the dynamo type. Direct current may flow in a conductor such as a wire, but can also be through semiconductors, insulators, or even through a vacuum as in electron or ion beams. The electric charge flows in a constant direction, distinguishing it from alternating current (AC). A term formerly used for direct current was Galvanic current.Direct current may be obtained from an alternating current supply by use of a current-switching arrangement called a rectifier, which contains electronic elements (usually) or electromechanical elements (historically) that allow current to flow only in one direction. Direct current may be made into alternating current with an inverter or a motor-generator set.
The first commercial electric power transmission (developed by Thomas Edison in the late nineteenth century) used direct current. Because of the advantage of alternating current over direct current in transforming and transmission, electric power distribution today is nearly all alternating current. For applications requiring direct current, such as third rail power systems, alternating current is distributed to a substation, which utilizes a rectifier to convert the power to direct current. See War of Currents. Direct current is used to charge batteries, and in nearly all electronic systems as the power supply. Very large quantities of direct-current power are used in production of aluminum and other electrochemical processes. Direct current is used for some railway propulsion, especially in urban areas. High voltage direct current is used to transmit large amounts of power from remote generation sites or to interconnect alternating current power grids.
Applications:-

Direct-current installations usually have different types of sockets, switches, and fixtures, mostly due to the low voltages used, from those suitable for alternating current. It is usually important with a direct-current appliance not to reverse polarity unless the device has a diode bridge to correct for this (most battery-powered devices do not).DC is commonly found in all low-voltage applications, especially where these are powered by batteries, which can produce only DC, or solar power systems. Most automotive applications use DC, although the alternator is an AC device which uses a rectifier to produce DC. Most electronic circuits require a DC power supply. Applications using fuel cells (mixing hydrogen and oxygen together with a catalyst to produce electricity and water as byproducts) also produce only DC.
Many telephones connect to a twisted pair of wires, and internally separate the AC component of the voltage between the two wires (the audio signal) from the DC component of the voltage between the two wires (used to power the phone).Telephone exchange communication equipment, such as DSLAM, uses standard -48V DC power supply. The negative polarity is achieved by grounding the positive terminal of power supply system and the battery bank. This is done to prevent electrolysis depositions.An electrified third rail can be used to power both underground (subway) and overground trains.

Alternating current

2009-11-20T01:04:00.000-08:00

In alternating current (AC, also ac) the movement (or flow) of electric charge periodically reverses direction. An electric charge would for instance move forward, then backward, then forward, then backward, over and over again. In direct current (DC), the movement (or flow) of electric charge is only in one direction.Used generically, AC refers to the form in which electricity is delivered to businesses and residences. The usual waveform of an AC power circuit is a sine wave, however in certain applications, different waveforms are used, such as triangular or square waves. Audio and radio signals carried on electrical wires are also examples of alternating current. In these applications, an important goal is often the recovery of information encoded (or modulated) onto the AC signal.

History:-

A power transformer developed by Lucien Gaulard and John Dixon Gibbs was demonstrated in London in 1881, and attracted the interest of Westinghouse. They also exhibited the invention in Turin in 1884, where it was adopted for an electric lighting system. Many of their designs were adapted to the particular laws governing electrical distribution in the UK.[citation needed]
In 1882, 1884, and 1885 Gaulard and Gibbs applied for patents on their transformer; however, these were overturned due to prior arts of Nikola Tesla and actions initiated by Sebastian Ziani de Ferranti.Ferranti went into this business in 1882 when he set up a shop in London designing various electrical devices. Ferranti believed in the success of alternating current power distribution early on, and was one of the few experts in this system in the UK. In 1887 the London Electric Supply Corporation (LESCo) hired Ferranti for the design of their power station at Deptford. He designed the building, the generating plant and the distribution system. On its completion in 1891 it was the first truly modern power station, supplying high-voltage AC power that was then "stepped down" for consumer use on each street. This basic system remains in use today around the world. Many homes all over the world still have electric meters with the Ferranti AC patent stamped on them.
William Stanley, Jr. designed one of the first practical devices to transfer AC power efficiently between isolated circuits. Using pairs of coils wound on a common iron core, his design, called an induction coil, was an early transformer. The AC power system used today developed rapidly after 1886, and includes key concepts by Nikola Tesla, who subsequently sold his patent to George Westinghouse. Lucien Gaulard, John Dixon Gibbs, Carl Wilhelm Siemens and others contributed subsequently to this field. AC systems overcame the limitations of the direct current system used by Thomas Edison to distribute electricity efficiently over long distances even though Edison attempted to discredit alternating current as too dangerous during the War of Currents.The first commercial power plant in the United States using three-phase alternating current was at the Mill Creek hydroelectric plant near Redlands, California in 1893 designed by Almirian Decker. Decker's design incorporated 10,000-volt three-phase transmission and established the standards for the complete system of generation, transmission and motors used today.The Jaruga power plant in Croatia was set in operation on 28 August 1895, . It was completed three days after the Niagara Falls plant, becoming the second commercial hydro power plant in the world. The two generators (42 Hz, 550 kW each) and the transformers were produced and installed by the Hungarian company Ganz. The transmission line from the power plant to the City of Šibenik was 11.5 kilometers (7.1 mi) long on wooden towers, and the municipal distribution grid 3000 V/110 V included six transforming stations. Alternating current circuit theory evolved rapidly in the latter part of the 19th and early 20th century. Notable contributors to the theoretical basis of alternating current calculations include Charles Steinmetz, James Clerk Maxwell, Oliver Heaviside, and many others. Calculations in unbalanced three-phase systems were simplified by the symmetrical components methods discussed by Charles Legeyt Fortescue in 1918.

Light-emitting diode

2009-11-18T02:16:00.000-08:00

A light-emitting diode (LED) (pronounced /ˌɛl.iːˈdiː/[1], or just /lɛd/), is an electronic light source. LEDs are used as indicator lamps in many kinds of electronics and increasingly for lighting. LEDs work by the effect of electroluminescence, discovered by accident in 1907. The LED was introduced as a practical electronic component in 1962.All early devices emitted low-intensity red light, but modern LEDs are available across the visible, ultraviolet and infra red wavelengths, with very high brightness.
LEDs are based on the semiconductor diode. When the diode is forward biased (switched on), electrons are able to recombine with holes and energy is released in the form of light. This effect is called electroluminescence and the color of the light is determined by the energy gap of the semiconductor. The LED is usually small in area (less than 1 mm2) with integrated optical components to shape its radiation pattern and assist in reflection.
LEDs present many advantages over traditional light sources including lower energy consumption, longer lifetime, improved robustness, smaller size and faster switching. However, they are relatively expensive and require more precise current and heat management than traditional light sources.
Applications of LEDs are diverse. They are used as low-energy indicators but also for replacements for traditional light sources in general lighting, automotive lighting and traffic signals. The compact size of LEDs has allowed new text and video displays and sensors to be developed, while their high switching rates are useful in communications technology.
History:-
Electroluminescence was discovered in 1907 by the British experimenter H. J. Round of Marconi Labs, using a crystal of silicon carbide and a cat's-whisker detector. Russian Oleg Vladimirovich Losev independently reported on the creation of a LED in 1927. His research was distributed in Russian, German and British scientific journals, but no practical use was made of the discovery for several decades. Rubin Braunstein of the Radio Corporation of America reported on infrared emission from gallium arsenide (GaAs) and other semiconductor alloys in 1955. Braunstein observed infrared emission generated by simple diode structures using gallium antimonide (GaSb), GaAs, indium phosphide (InP), and silicon-germanium (SiGe) alloys at room temperature and at 77 kelvin.
In 1961, experimenters Robert Biard and Gary Pittman working at Texas Instrumentsfound that GaAs emitted infrared radiation when electric current was applied and received the patent for the infrared LED.
The first practical visible-spectrum (red) LED was developed in 1962 by Nick Holonyak Jr., while working at General Electric Company.[2] Holonyak is seen as the "father of the light-emitting diode". M. George Craford a former graduate student of Holonyak, invented the first yellow LED and improved the brightness of red and red-orange LEDs by a factor of ten in 1972.[13] In 1976, T.P. Pearsall created the first high-brightness, high efficiency LEDs for optical fiber telecommunications by inventing new semiconductor materials specifically adapted to optical fiber transmission wavelengths.
Up to 1968 visible and infrared LEDs were extremely costly, on the order of US $200 per unit, and so had little practical application.The Monsanto Company was the first organization to mass-produce visible LEDs, using gallium arsenide phosphide in 1968 to produce red LEDs suitable for indicators. Hewlett Packard (HP) introduced LEDs in 1968, initially using GaAsP supplied by Monsanto. The technology proved to have major applications for alphanumeric displays and was integrated into HP's early handheld calculators.
The first commercial LEDs were commonly used as replacements for incandescent indicators, and in seven-segment displays, first in expensive equipment such as laboratory and electronics test equipment, then later in such appliances as TVs, radios, telephones, calculators, and even watches (see list of signal applications). These red LEDs were bright enough only for use as indicators, as the light output was not enough to illuminate an area. Later, other colors became widely available and also appeared in appliances and equipment. As the LED materials technology became more advanced, the light output was increased, while maintaining the efficiency and the reliability to an acceptable level. The invention and development of the high power white light LED led to use for illumination(see list of illumination applications). Most LEDs were made in the very common 5 mm T1¾ and 3 mm T1 packages, but with increasing power output, it has become increasingly necessary to shed excess heat in order to maintain reliability[19], so more complex packages have been adapted for efficient heat dissipation. Packages for state-of-the-art high power LEDs bear little resemblance to early LEDs.
The first high-brightness blue LED was demonstrated by Shuji Nakamura of Nichia Corporation and was based on InGaN borrowing on critical developments in GaN nucleation on sapphire substrates and the demonstration of p-type doping of GaN which were developed by Isamu Akasaki and H. Amano in Nagoya. In 1995, Alberto Barbieri at the Cardiff University Laboratory (GB) investigated the efficiency and reliability of high-brightness LEDs and demonstrated a very impressive result by using a transparent contact made of indium tin oxide (ITO) on (AlGaInP/GaAs) LED. The existence of blue LEDs and high efficiency LEDs quickly led to the development of the first white LED, which employed a Y3Al5O12:Ce, or "YAG", phosphor coating to mix yellow (down-converted) light with blue to produce light that appears white. Nakamura was awarded the 2006 Millennium Technology Prize for his invention.
The development of LED technology has caused their efficiency and light output to increase exponentially, with a doubling occurring about every 36 months since the 1960s, in a way similar to Moore's law. The advances are generally attributed to the parallel development of other semiconductor technologies and advances in optics and material science. This trend is normally called Haitz's Law after Dr. Roland Haitz. In February 2008, Bilkent university in Turkey reported 300 lumens of visible light per watt luminous efficacy (not per electrical watt) and warm light by using nanocrystals.
In January 2009, researchers from Cambridge University reported a process for growing gallium nitride (GaN) LEDs on silicon. Production costs could be reduced by 90% using six-inch silicon wafers instead of two-inch sapphire wafers. The team was led by Colin Humphreys.

Galvanometer

2009-11-18T01:52:00.000-08:00

A galvanometer is a type of ammeter: an instrument for detecting and measuring electric
current. It is an analog electromechanical transducer that produces a rotary deflection of some type of pointer in response to electric current flowing through its coil. The term has expanded to include uses of the same mechanism in recording, positioning, and servomechanism equipment. History:-
The deflection of a magnetic compass needle by current in a wire was first described by Hans Oersted in 1820. The phenomenon was studied both for its own sake and as a means of measuring electrical current. The earliest galvanometer was reported by Johann Schweigger at the University of Halle on 16 September 1820. André-Marie Ampère also contributed to its development. Early designs increased the effect of the magnetic field due to the current by using multiple turns of wire; the instruments were at first called "multipliers" due to this common design feature. The term "galvanometer", in common use by 1836, was derived from the surname of Italian electricity researcher Luigi Galvani, who discovered that electric current could make a frog's leg jerk.
Originally the instruments relied on the Earth's magnetic field to provide the restoring force for the compass needle; these were called "tangent" galvanometers and had to be oriented before use. Later instruments of the "astatic" type used opposing magnets to become independent of the Earth's field and would operate in any orientation. The most sensitive form, the Thompson or mirror galvanometer, was invented by William Thomson (Lord Kelvin). Instead of a compass needle, it used tiny magnets attached to a small lightweight mirror, suspended by a thread; the deflection of a beam of light greatly magnified the deflection due to small currents. Alternatively the deflection of the suspended magnets could be observed directly through a microscope.
The ability to quantitatively measure voltage and current allowed Georg Ohm to formulate Ohm's Law, which states that the voltage across an element is directly proportional to the current through it.
The early moving-magnet form of galvanometer had the disadvantage that it was affected by any magnets or iron masses near it, and its deflection was not linearly proportional to the current. In 1882 Jacques-Arsène d'Arsonval developed a form with a stationary permanent magnet and a moving coil of wire, suspended by coiled hair springs. The concentrated magnetic field and delicate suspension made these instruments sensitive and they could be mounted in any position. By 1888 Edward Weston had brought out a commercial form of this instrument, which became a standard component in electrical equipment. This design is almost universally used in moving-vane meters today.
Opration:-
The most familiar use is as an analog measuring instrument, often called a meter. It is used to measure the direct current (flow of electric charge) through an electric circuit. The D'Arsonval/Weston form used today is constructed with a small pivoting coil of wire in the field of a permanent magnet. The coil is attached to a thin pointer that traverses a calibrated scale. A tiny torsion spring pulls the coil and pointer to the zero position.
When a direct current (DC) flows through the coil, the coil generates a magnetic field. This field acts against the permanent magnet. The coil twists, pushing against the spring, and moves the pointer. The hand points at a scale indicating the electric current. Careful design of the pole pieces ensures that the magnetic field is uniform, so that the angular deflection of the pointer is proportional to the current. A useful meter generally contains provision for damping the mechanical resonance of the moving coil and pointer, so that the pointer settles quickly to its position without oscillation.
The basic sensitivity of a meter might be, for instance, 100 microamperes full scale (with a voltage drop of, say, 50 millivolts at full current). Such meters are often calibrated to read some other quantity that can be converted to a current of that magnitude. The use of current dividers, often called shunts, allows a meter to be calibrated to measure larger currents. A meter can be calibrated as a DC voltmeter if the resistance of the coil is known by calculating the voltage required to generate a full scale current. A meter can be configured to read other voltages by putting it in a voltage divider circuit. This is generally done by placing a resistor in series with the meter coil. A meter can be used to read resistance by placing it in series with a known voltage (a battery) and an adjustable resistor. In a preparatory step, the circuit is completed and the resistor adjusted to produce full scale deflection. When an unknown resistor is placed in series in the circuit the current will be less than full scale and an appropriately calibrated scale can display the value of the previously-unknown resistor.
Because the pointer of the meter is usually a small distance above the scale of the meter, parallax error can occur when the operator attempts to read the scale line that "lines up" with the pointer. To counter this, some meters include a mirror along the markings of the principal scale. The accuracy of the reading from a mirrored scale is improved by positioning one's head while reading the scale so that the pointer and the reflection of the pointer are aligned; at this point, the operator's eye must be directly above the pointer and any parallax error has been minimized.
Types:-
Extremely sensitive measuring equipment once used mirror galvanometers that substituted a mirror for the pointer. A beam of light reflected from the mirror acted as a long, massless pointer. Such instruments were used as receivers for early trans-Atlantic telegraph systems, for instance. The moving beam of light could also be used to make a record on a moving photographic film, producing a graph of current versus time, in a device called an oscillograph.
Today the main type of galvanometer mechanism still used is the moving coil D'Arsonval/Weston mechanism, which is used in traditional analog meters.
Tangent galvanometer:-
A tangent galvanometer is an early measuring instrument used for the measurement of electric current. It works by using a compass needle to compare a magnetic field generated by the unknown current to the magnetic field of the Earth. It gets its name from its operating principle, the tangent law of magnetism, which states that the tangent of the angle a compass needle makes is proportional to the ratio of the strengths of the two perpendicular magnetic fields. It was first described by Claude Pouillet in 1837.
A tangent galvanometer consists of a coil of insulated copper wire wound on a circular non-magnetic frame. The frame is mounted vertically on a horizontal base provided with levelling screws. The coil can be rotated on a vertical axis passing through its centre. A compass box is mounted horizontally at the centre of a circular scale. It consists of a tiny, powerful magnetic needle pivoted at the centre of the coil. The magnetic needle is free to rotate in the horizontal plane. The circular scale is divided into four quadrants. Each quadrant is graduated from 0° to 90°. A long thin aluminium pointer is attached to the needle at its centre and at right angle to it. To avoid errors due to parallax a plane mirror is mounted below the compass needle.
In operation, the instrument is first rotated until the magnetic field of the Earth, indicated by the compass needle, is parallel with the plane of the coil. Then the unknown current is applied to the coil. This creates a second magnetic field on the axis of the coil, perpendicular to the Earth's magnetic field. The compass needle responds to the vector sum of the two fields, and deflects to an angle equal to the tangent of the ratio of the two fields. From the angle read from the compass's scale, the current could be found from a table.
The current supply wires have to be wound in a small helix, like a pig's tail, otherwise the field due to the wire will affect the compass needle and an incorrect reading will be obtained.
uses:-
A major early use for galvanometers was for finding faults in telecommunications cables. They were superseded in this application late in the 20th century by time-domain reflectometers.
Probably the largest use of galvanometers was the D'Arsonval/Weston type movement used in analog meters in electronic equipment. Since the 1980s, galvanometer-type analog meter movements have been displaced by analog to digital converters (ADCs) for some uses. A digital panel meter (DPM) contains an analog to digital converter and numeric display. The advantages of a digital instrument are higher precision and accuracy, but factors such as power consumption or cost may still favor application of analog meter movements.

Three-phase electric power

2009-11-15T08:57:00.000-08:00

Three-phase electric power is a common method of alternating-current electric power transmission.It is a type of polyphase system, and is the most common method used by electric power distribution grids worldwide to distribute power. It is also used to power large motors and other large loads. A three-phase system is generally more economical than others because it uses less conductor material to transmit electric power than equivalent single-phase or two-phase systems at the same voltage.
In a three-phase system, three circuit conductors carry three alternating currents (of the same frequency) which reach their instantaneous peak values at different times. Taking one conductor as the reference, the other two currents are delayed in time by one-third and two-thirds of one cycle of the electrical current. This delay between phases has the effect of giving constant power transfer over each cycle of the current, and also makes it possible to produce a rotating magnetic field in an electric motor.Three-phase systems may or may not have a neutral wire. A neutral wire allows the three-phase system to use a higher voltage while still supporting lower-voltage single-phase appliances. In high-voltage distribution situations, it is common not to have a neutral wire as the loads can simply be connected between phases (phase-phase connection).
Three-phase has properties that make it very desirable in electric power systems:
The phase currents tend to cancel out one another, summing to zero in the case of a linear balanced load. This makes it possible to eliminate or reduce the size of the neutral conductor; all the phase conductors carry the same current and so can be the same size, for a balanced load. Power transfer into a linear balanced load is constant, which helps to reduce generator and motor vibrations. Three-phase systems can produce a magnetic field that rotates in a specified direction, which simplifies the design of electric motors. Three is the lowest phase order to exhibit all of these properties.Most household loads are single-phase. In North America and some other countries, three-phase power generally does not enter homes. Even in areas where it does, it is typically split out at the main distribution board and the individual loads are fed from a single phase. Sometimes it is used to power electric stoves and washing machines.The three phases are typically indicated by colors which vary by country. See the table for more information

Two-phase electrical power

2009-11-15T08:54:00.000-08:00

Two-phase electrical power was an early 20th century polyphase alternating current electric power distribution system. Two circuits were used, with voltage phases differing by 90 degrees. Usually circuits used four wires, two for each phase. Less frequently, three wires were used, with a common wire with a larger-diameter conductor. Some early two-phase generators had two complete rotor and field assemblies, with windings physically offset by 90 electrical degrees to provide two-phase power. The generators at Niagara Falls installed in 1895 were the largest generators in the world at the time and were two-phase machines.
The advantage of two-phase electrical power was that it allowed for simple, self-starting electric motors. In the early days of electrical engineering, it was easier to analyze and design two-phase systems where the phases were completely separated. [1] It was not until the invention of the method of symmetrical components in 1918 that polyphase power systems had a convenient mathematical tool for describing unbalanced load cases. The revolving magnetic field produced with a two-phase system allowed electric motors to provide torque from zero motor speed, which was not possible with a single-phase induction motor (without extra starting means). Induction motors designed for two-phase operation use the same winding configuration as capacitor start single-phase motors.
Three-phase electric power requires less conductor mass for the same voltage and overall amount of power, compared with a two-phase four-wire circuit of the same carrying capacity. It has all but replaced two-phase power for commercial distribution of electrical energy, but two-phase circuits are still found in certain control systems.
Two-phase circuits typically use two separate pairs of current-carrying conductors. Alternatively, three wires may be used, but the common conductor carries the vector sum of the phase currents, which requires a larger conductor. Three phase can share conductors so that the three phases can be carried on three conductors of the same size. In electrical power distribution, a requirement of only three conductors rather than four represented a considerable distribution-wire cost savings due to the expense of conductors and installation.
Two-phase power can be derived from a three-phase source using two transformers in a Scott connection. One transformer primary is connected across two phases of the supply. The second transformer is connected to a center-tap of the first transformer, and is wound for 86.6% of the phase-to-phase voltage on the 3-phase system. The secondaries of the transformers will have two phases 90 degrees apart in time, and a balanced two-phase load will be evenly balanced over the three supply phases.Three-wire, 120/240 volt single phase power used in the USA and Canada is sometimes incorrectly called "two-phase". The proper term is split phase or 3-wire single-phase.

Single-phase electric power

2009-11-15T08:49:00.001-08:00

In electrical engineering, single-phase electric power refers to the distribution of alternating current electric power using a system in which all the voltages of the supply vary in unison. Single-phase distribution is used when loads are mostly lighting and heating, with few large electric motors. A single-phase supply connected to an alternating current electric motor does not produce a revolving magnetic field; single-phase motors need additional circuits for starting, and such motors are uncommon above 10 or 20 kW in rating.
In contrast, in a three-phase system, the currents in each conductor reach their peak instantaneous values sequentially, not simultaneously; in each cycle of the power frequency, first one, then the second, then the third current reaches its maximum value. The waveforms of the three supply conductors are offset from one another in time (delayed in phase) by one-third of their period.Standard frequencies of single-phase power systems are either 50 or 60 Hz. Special single-phase traction power networks may operate at 16.67 Hz or other frequencies to power electric railways.
Splitting out:-

No arrangement of transformers can convert a single-phase load into a balanced load on a polyphase system. A single-phase load may be powered from a three-phase distribution system either by connection between a phase and neutral or by connecting the load between two phases. The load device must be designed for the voltage in each case. The neutral point in a three phase system exists at the mathematical center of an equilateral triangle formed by the three phase points, and the phase-to-phase voltage is accordingly times the phase-to-neutral voltage.[1] For example, in places using a 415 volt 3 phase system, the phase-to-neutral voltage is 240 volts, allowing single-phase lighting to be connected phase-to-neutral and three-phase motors to be connected to all three phases.
In North America, a typical three-phase system will have 208 volts between the phases and 120 volts between phase and neutral. If heating equipment designed for the 240-volt three-wire single phase system is connected to two phases of a 208 volt supply, it will only produce 75% of its rated heating effect. Single-phase motors may have taps to allow their use on either 208 V or 240 V supplies.
On higher voltage systems (kilovolts) where a single phase transformer is in use to supply a low voltage system the method of splitting varies. In North America utility distribution practice, the primary of the step-down transformer is wired across a single high voltage feed wire and neutral, at least for smaller supplies (see photo of transformer on right). Rural distribution may be a single phase at a medium voltage; in some areas single wire earth return distribution is used when customers are very far apart. In Britain the step-down primary is wired phase-phase.
Applications:-

Single-phase power distribution is widely used especially in rural areas, where the cost of a three-phase distribution network is high and motor loads are small and uncommon.
High power systems, say, hundreds of kVA or larger, are nearly always three phase. The largest supply normally available as single phase varies according to the standards of the electrical utility. In the UK a single-phase household supply may be rated 100 A or even 125 A, meaning that there is little need for 3 phase in a domestic or small commercial environment. Much of the rest of Europe has traditionally had much smaller limits on the size of single phase supplies resulting in even houses being supplied with 3 phase (in urban areas with three-phase supply networks).
In North America, individual residences and small commercial buildings with services up to about 100 kV·A (417 amperes at 240 volts) will usually have three-wire single-phase distribution, often with only one customer per distribution transformer. In exceptional cases larger single-phase three-wire services can be provided, usually only in remote areas where polyphase distribution is not available. In rural areas farmers who wish to use three-phase motors may install a phase converter if only a single-phase supply is available. Larger consumers such as large buildings, shopping centres, factories, office blocks, and multiple-unit apartment blocks will have three-phase service. In densely-populated areas of cities, network power distribution is used with many customers and many supply transformers connected to provide hundreds or thousands of kV·A load concentrated over a few hundred square metres.
Three-wire single-phase systems are rarely used in the UK where large loads are needed off only two high voltage phases.
Single-phase power may be used for electric railways; the largest single-phase generator in the world, at Neckarwestheim Nuclear Power Plant, supplies a railway system on a dedicated traction power network.

Electric field

2009-11-15T08:43:00.000-08:00

In physics, the space surrounding an electric charge or in the presence of a time-varying magnetic field has a property called an electric field. This electric field exerts a force on other electrically charged objects. The concept of an electric field was introduced by Michael Faraday.
The electric field is a vector field with SI units of newtons per coulomb (N C−1) or, equivalently, volts per metre (V m−1). The SI base units of the electric field are kg·m·s−3·A−1. The strength of the field at a given point is defined as the force that would be exerted on a positive test charge of +1 coulomb placed at that point; the direction of the field is given by the direction of that force. Electric fields contain electrical energy with energy density proportional to the square of the field intensity. The electric field is to charge as gravitational acceleration is to mass and force density is to volume.A moving charge has not just an electric field but also a magnetic field, and in general the electric and magnetic fields are not completely separate phenomena; what one observer perceives as an electric field, another observer in a different frame of reference perceives as a mixture of electric and magnetic fields. For this reason, one speaks of "electromagnetism" or "electromagnetic fields." In quantum mechanics, disturbances in the electromagnetic fields are called photons, and the energy of photons is quantized.

Wattmeter

2009-11-10T09:41:00.000-08:00

The wattmeter is an instrument for measuring the electric power (or the supply rate of electrical energy) in watts of any given circuit.
Electrodynamic:-
The traditional analog wattmeter is an electrodynamic instrument. The device consists of a pair of fixed coils, known as current coils, and a movable coil known as the potential coil.
The current coils connected in series with the circuit, while the potential coil is connected in parallel. Also, on analog wattmeters, the potential coil carries a needle that moves over a scale to indicate the measurement. A current flowing through the current coil generates an electromagnetic field around the coil. The strength of this field is proportional to the line current and in phase with it. The potential coil has, as a general rule, a high-value resistor connected in series with it to reduce the current that flows through it.
The result of this arrangement is that on a dc circuit, the deflection of the needle is proportional to both the current and the voltage, thus conforming to the equation W=VA or P=VI. On an ac circuit the deflection is proportional to the average instantaneous product of voltage and current, thus measuring true power, and possibly (depending on load characteristics) showing a different reading to that obtained by simply multiplying the readings showing on a stand-alone voltmeter and a stand-alone ammeter in the same circuit.
The two circuits of a wattmeter can be damaged by excessive current. The ammeter and voltmeter are both vulnerable to overheating — in case of an overload, their pointers will be driven off scale — but in the wattmeter, either or even both the current and potential circuits can overheat without the pointer approaching the end of the scale! This is because the position of the pointer depends on the power factor, voltage and current. Thus, a circuit with a low power factor will give a low reading on the wattmeter, even when both of its circuits are loaded to the maximum safety limit. Therefore, a wattmeter is rated not only in watts, but also in volts and amperes.
Electrodynamometer:-
An early current meter was the electrodynamometer. Used in the early 20th century, the Siemens electrodynamometer, for example, is a form of an electrodynamic ammeter, that has a fixed coil which is surrounded by another coil having its axis at right angles to that of the fixed coil. This second coil is suspended by a number of silk fibres, and to the coil is also attached a spiral spring the other end of which is fastened to a torsion head. If then the torsion head is twisted, the suspended coil experiences a torque and is displaced through and angle equal to that of the torsion head. The current can be passed into and out of the movable coil by permitting the ends of the coil to dip into two mercury cups.
If a current is passed through the fixed coil and movable coil in series with one another, the movable coil tends to displace itself so as to bring the axes of the coils, which are normally at right angles, more into the same direction. This tendency can be resisted by giving a twist to the torsion head and so applying to the movable coil through the spring a restoring torque, which opposes the torque due to the dynamic action of the currents. If then the torsion head is provided with an index needle, and also if the movable coil is provided with an indicating point, it is possible to measure the torsional angle through which the head must be twisted to bring the movable coil back to its zero position. In these circumstances, the torsional angle becomes a measure of the torque and therefore of the product of the strengths of the currents in the two coils, that is to say, of the square of the strength of the current passing through the two coils if they are joined up in series. The instrument can therefore be graduated by passing through it known and measured continuous currents, and it then becomes available for use with either continuous or alternating currents. The instrument can be provided with a curve or table showing the current corresponding to each angular displacement of the torsion head.

Integrated circuit

2009-11-10T09:36:00.000-08:00

In electronics, an integrated circuit (also known as IC, microcircuit, microchip, silicon chip, or chip) is a miniaturized electronic circuit (consisting mainly of semiconductor devices, as well as passive components) that has been manufactured in the surface of a thin substrate of semiconductor material. Integrated circuits are used in almost all electronic equipment in use today and have revolutionized the world of electronics.
A hybrid integrated circuit is a miniaturized electronic circuit constructed of individual semiconductor devices, as well as passive components, bonded to a substrate or circuit board.
Introduction:-
Integrated circuits were made possible by experimental discoveries which showed that semiconductor devices could perform the functions of vacuum tubes, and by mid-20th-century technology advancements in semiconductor device fabrication. The integration of large numbers of tiny transistors into a small chip was an enormous improvement over the manual assembly of circuits using discrete electronic components. The integrated circuit's mass production capability, reliability, and building-block approach to circuit design ensured the rapid adoption of standardized ICs in place of designs using discrete transistors.
There are two main advantages of ICs over discrete circuits: cost and performance. Cost is low because the chips, with all their components, are printed as a unit by photolithography and not constructed one transistor at a time. Furthermore, much less material is used to construct a circuit as a packaged IC die than as a discrete circuit. Performance is high since the components switch quickly and consume little power (compared to their discrete counterparts) because the components are small and close together. As of 2006, chip areas range from a few square millimeters to around 350 mm2, with up to 1 million transistors per mm2.
Invantion:-
The idea of an integrated circuit was conceived by a radar scientist working for the Royal Radar Establishment of the British Ministry of Defence, Geoffrey W.A. Dummer (1909-2002), who published it at the Symposium on Progress in Quality Electronic Components in Washington, D.C. on May 7, 1952. He gave many symposia publicly to propagate his ideas.
Dummer unsuccessfully attempted to build such a circuit in 1956.
The integrated circuit can be credited as being invented by both Jack Kilby of Texas Instruments and Robert Noyce of Fairchild Semiconductor [3] working independently of each other. Kilby recorded his initial ideas concerning the integrated circuit in July 1958 and successfully demonstrated the first working integrated circuit on September 12, 1958. In his patent application of February 6, 1959, Kilby described his new device as “a body of semiconductor material ... wherein all the components of the electronic circuit are completely integrated.”
Kilby won the 2000 Nobel Prize in Physics for his part of the invention of the integrated circuit. Robert Noyce also came up with his own idea of integrated circuit, half a year later than Kilby. Noyce's chip had solved many practical problems that the microchip developed by Kilby had not. Noyce's chip, made at Fairchild, was made of silicon, whereas Kilby's chip was made of germanium.
Early developments of the integrated circuit go back to 1949, when the German engineer Werner Jacobi (Siemens AG) filed a patent for an integrated-circuit-like semiconductor amplifying device showing five transistors on a common substrate arranged in a 2-stage amplifier arrangement. Jacobi discloses small and cheap hearing aids as typical industrial applications of his patent. A commercial use of his patent has not been reported.
A precursor idea to the IC was to create small ceramic squares (wafers), each one containing a single miniaturized component. Components could then be integrated and wired into a bidimensional or tridimensional compact grid. This idea, which looked very promising in 1957, was proposed to the US Army by Jack Kilby, and led to the short-lived Micromodule Program (similar to 1951's Project Tinkertoy). However, as the project was gaining momentum, Kilby came up with a new, revolutionary design: the IC.
The aforementioned Noyce credited Kurt Lehovec of Sprague Electric for the principle of p-n junction isolation caused by the action of a biased p-n junction (the diode) as a key concept behind the IC.

History of electrical engineering

2009-11-10T09:31:00.000-08:00

This article details the history of electrical engineering. Topics also included are the general developments and notable individuals within the electrical engineering profession.
History:-

Thales of Miletus, an ancient greek philosopher, writing at around 600 BCE, described a form of static electricity, noting that rubbing fur on various substances, such as amber, would cause a particular attraction between the two. He noted that the amber buttons could attract light objects such as hair and that if they rubbed the amber for long enough they could even get a spark to jump.At around 450 B.C. Democritus, a later Greek philosopher, developed an atomic theory that was remarkably similar to our modern atomic theory. His mentor, Leucippus, is credited with this same theory. The hypothesis of Leucippus and Democritus held everything to be composed of atoms. But these atoms, called "atomos", were indivisible, and indestructible. He presciently stated that between atoms lies empty space, and that atoms are constantly in motion. He was incorrect only in stating that atoms come different sizes and shapes. Each object had its own shaped and sized atom.
An object found in Iraq in 1938, dated to about 250 BCE and called the Baghdad Battery, resembles a galvanic cell and is believed by some to have been used for electroplating in Mesopotamia, although this has not yet been proven.
Thales of Miletus, an ancient greek philosopher, writing at around 600 BCE, described a form of static electricity, noting that rubbing fur on various substances, such as amber, would cause a particular attraction between the two. He noted that the amber buttons could attract light objects such as hair and that if they rubbed the amber for long enough they could even get a spark to jump.
At around 450 B.C. Democritus, a later Greek philosopher, developed an atomic theory that was remarkably similar to our modern atomic theory. His mentor, Leucippus, is credited with this same theory. The hypothesis of Leucippus and Democritus held everything to be composed of atoms. But these atoms, called "atomos", were indivisible, and indestructible. He presciently stated that between atoms lies empty space, and that atoms are constantly in motion. He was incorrect only in stating that atoms come different sizes and shapes. Each object had its own shaped and sized atom.
An object found in Iraq in 1938, dated to about 250 BCE and called the Baghdad Battery, resembles a galvanic cell and is believed by some to have been used for electroplating in Mesopotamia, although this has not yet been proven.
During the latter part of the 1800s, the study of electricity was largely considered to be a subfield of physics. It was not until the late 19th century that universities started to offer degrees in electrical engineering. In 1882, Darmstadt University of Technology founded the first chair and the first faculty of electrical engineering worldwide. In the same year, under Professor Charles Cross, the Massachusetts Institute of Technology began offering the first option of Electrical Engineering within a physics department. In 1883, Darmstadt University of Technology and Cornell University introduced the world's first courses of study in electrical engineering and in 1885 the University College London founded the first chair of electrical engineering in the United Kingdom. The University of Missouri subsequently established the first department of electrical engineering in the United States in 1886.
During this period work in the area increased dramatically. In 1882 Edison switched on the world's first large-scale electrical supply network that provided 110 volts direct current to fifty-nine customers in lower Manhattan. In 1887 Nikola Tesla filed a number of patents related to a competing form of power distribution known as alternating current. In the following years a bitter rivalry between Tesla and Edison, known as the "War of Currents", took place over the preferred method of distribution. AC eventually replaced DC for generation and power distribution, enormously extending the range and improving the safety and efficiency of power distribution.The efforts of the two did much to further electrical engineering—Tesla's work on induction motors and polyphase systems influenced the field for years to come, while Edison's work on telegraphy and his development of the stock ticker proved lucrative for his company, which ultimately became General Electric.However, by the end of the 19th century, other key figures in the progress of electrical engineering were beginning to emerge.Charles Proteus Steinmetz helped foster the development of alternating current that made possible the expansion of the electric power industry in the United States, formulating mathematical theories for engineers.

Hybrid electric vehicle

2009-11-09T08:50:00.000-08:00

A hybrid electric vehicle (HEV) combines a conventional internal combustion engine propulsion system with an electric propulsion system. The presence of the electric powertrain is intended to achieve either better fuel economy than a conventional vehicle, or better performance. A variety of types of HEV exist, and the degree to which they function as EVs varies as well. The most common form of HEV is the hybrid electric car, although hybrid electric trucks (pickups and tractors) also exist.
Modern HEVs make use of efficiency-improving technologies such as regenerative braking, which converts the vehicle's kinetic energy into battery-replenishing electric energy, rather than wasting it as heat energy as conventional brakes do. Some varieties of HEVs use their internal combustion engine to generate electricity by spinning an electrical generator (this combination is known as a motor-generator), to either recharge their batteries or to directly power the electric drive motors. Many HEVs reduce idle emissions by shutting down the ICE at idle and restarting it when needed; this is known as a start-stop system. A hybrid-electric produces less emissions from its ICE than a comparably-sized gasoline car, as an HEV's gasoline engine is usually smaller than a pure fossil-fuel vehicle, and if not used to directly drive the car, can be geared to run at maximum efficiency, further improving fuel economy.
The hybrid-electric vehicle did not become widely available until the release of the Toyota Prius in Japan in 1997, followed by the Honda Insight in 1999. While initially perceived as unnecessary due to the low cost of gasoline, worldwide increases in the price of petroleum caused many automakers to release hybrids in the late 2000s; they are now perceived as a core segment of the automotive market of the future.Worldwide sales of hybrid vehicles produced by Toyota reached 1.0 million vehicles by May 31, 2007, and the 2.0 million mark was reached by August 31, 2009, with hybrids sold in 50 countries. Worldwide sales are led by the Prius, with cumulative sales of 1.43 million by Augut 2009. The second-generation Honda Insight was the top-selling vehicle in Japan in April 2009, marking the first occasion that an HEV has received the distinction. American automakers have made development of hybrid cars a top priority.
History:-
In 1901, while employed at Lohner Coach Factory, Ferdinand Porsche designed the Mixte, a 4WD series-hybrid version of "System Lohner-Porsche" electric carriage previously appeared in 1900 Paris Salon. The Mixte included a pair of generators driven by 2.5-hp Daimler IC engines to extend operating range. The Mixte broke several Austrian speed records, and also won the Exelberg Rally in 1901 with Porsche himself driving. The Mixte used a gasoline engine powering a generator, which in turn powered electric hub motors, with a small battery pack for reliability. It had a range of 50 km, a top speed of 50 km/h and a power of 5.22 kW during 20 minutes.
In 1905, H. Piper filed a US patent application for a hybrid vehicle.
The 1915 Dual Power, made by the Woods Motor Vehicle electric car maker, had a four-cylinder ICE and an electric motor. Below 15 mph (25 km/h) the electric motor alone drove the vehicle, drawing power from a battery pack, and above this speed the "main" engine cut in to take the car up to its 35 mph (55 km/h) top speed. About 600 were made up to 1918.
The first gasoline-electric hybrid car was released by the Woods Motor Vehicle Company of Chicago in 1917. The hybrid was a commercial failure, proving to be too slow for its price, and too difficult to service.
In 1931 Erich Gaichen invented and drove from Altenburg to Berlin a 1/2 horsepower electric car containing features later incorporated into hybrid cars. Its maximum speed was 25 miles per hour (40 km/h), but it was licensed by the Motor Transport Office, taxed by the German Revenue Department and patented by the German Reichs-Patent Amt. The car battery was re-charged by the motor when the car went downhill. Additional power to charge the battery was provided by a cylinder of compressed air which was re-charged by small air pumps activated by vibrations of the chassis and the brakes and by igniting oxyhydrogen gas. An account of the car and his characterization as a "crank inventor" can be found in Arthur Koestler's autobiography, Arrow in the Blue, pages 269-271, which summarize a contemporaneous newspaper account written by Koestler. No production beyond the prototype was reported.
Current technologey:-
A more recent working prototype of the HEV was built by Victor Wouk (one of the scientists involved with the Henney Kilowatt, the first transistor-based electric car). Wouk's work with HEVs in the 1960s and 1970s earned him the title as the "Godfather of the Hybrid".Wouk installed a prototype hybrid drivetrain (with a 16 kW electric motor) into a 1972 Buick Skylark provided by GM for the 1970 Federal Clean Car Incentive Program, but the program was stopped by the United States Environmental Protection Agency (EPA) in 1976 while Eric Stork, the head of the EPA at the time, was accused of a prejudicial coverup.
The regenerative braking system, the core design concept of most production HEVs, was developed by electrical engineer David Arthurs around 1978 using off-the shelf components and an Opel GT. However the voltage controller to link the batteries, motor (a jet-engine starter motor), and DC generator was Arthurs'. The vehicle exhibited 75 miles per US gallon (3.1 L/100 km; 90 mpg-imp) fuel efficiency and plans for it (as well as somewhat updated versions) are still available through the Mother Earth News web site. The Mother Earth News' own 1980 version claimed nearly 84 miles per US gallon (2.8 L/100 km; 101 mpg-imp).
In 1989, Audi produced its first iteration of the Audi Duo (or Audi 100 Avant duo) experimental vehicle, a plug-in parallel hybrid based on the Audi 100 Avant quattro. This car had a 12.6 bhp Siemens electric motor which drove the rear wheels. A trunk-mounted nickel-cadmium battery supplied energy to the motor that drove the rear wheels. The vehicle's front wheels were powered by a 2.3-litre five-cylinder engine with an output of 136 bhp (101 kW). The intent was to produce a vehicle which could operate on the engine in the country and electric mode in the city. Mode of operation could be selected by the driver. Just ten vehicles are believed to have been made; one drawback was that due to the extra weight of the electric drive, the vehicles were less efficient when running on their engines alone than standard Audi 100s with the same engine.
Two years later, Audi, unveiled the second duo generation - likewise based on the Audi 100 Avant quattro. Once again this featured an electric motor, a 28.6 bhp (21.3 kW) three-phase machine, driving the rear wheels. This time, however, the rear wheels were additionally powered via the Torsen differential from the main engine compartment, which housed a 2.0-litre four-cylinder engine.
The Bill Clinton administration initiated the Partnership for a New Generation of Vehicles (PNGV) program on 29 September 1993 that involved Chrysler, Ford, General Motors, USCAR, the DoE, and other various governmental agencies to engineer the next efficient and clean vehicle. The NRC cited automakers’ moves to produce HEVs as evidence that technologies developed under PNGV were being rapidly adopted on production lines, as called for under Goal 2. Based on information received from automakers, NRC reviewers questioned whether the “Big Three” would be able to move from the concept phase to cost effective, pre-production prototype vehicles by 2004, as set out in Goal 3. The program was replaced by the hydrogen-focused FreedomCAR initiative by the George W. Bush administration in 2001, an initiative to fund research too risky for the private sector to engage in, with the long-term goal of developing effectively carbon emission- and petroleum-free vehicles.

Capacitance

2009-11-09T08:11:00.000-08:00

In electromagnetism and electronics, capacitance is the ability of a body to hold an electrical charge. Capacitance is also a measure of the amount of electric charge stored (or separated) for a given electric potential. A common form of charge storage device is a parallel-plate capacitor. In a parallel plate capacitor capacitance is directly proportional to the surface area of the conductor plates and inversely proportional to the separation distance between the plates. If the charges on the plates are +Q and −Q, and V gives the voltage between the plates, then the capacitance is given by.

C=Q/V

Asynchronous circuit

2009-11-08T09:41:00.000-08:00

An asynchronous circuit is a circuit in which the parts are largely autonomous. They are not governed by a clock circuit or global clock signal, but instead need only wait for the signals that indicate completion of instructions and operations. These signals are specified by simple data transfer protocols. This digital logic design is contrasted with a synchronous circuit which operates according to clock timing signals.

Theoretical foundations:-
Petri Nets are an attractive and powerful model for reasoning about asynchronous circuits. However Petri nets have been criticized by Carl Hewitt for their lack of physical realism (see Petri net). Subsequent to Petri nets other models of concurrency have been developed that can model asynchronous circuits including the Actor model and process calculi.
The term asynchronous logic is used to describe a variety of design styles, which use different assumptions about circuit properties. These vary from the bundled delay model - which uses 'conventional' data processing elements with completion indicated by a locally generated delay model - to delay-insensitive design - where arbitrary delays through circuit elements can be accommodated. The latter style tends to yield circuits which are larger than bundled data implementations, but which are insensitive to layout and parametric variations and are thus "correct by design."

Benefits:-

Different classes of asynchronous circuitry offer different advantages. Below is a list of the advantages offered by Quasi Delay Insensitive Circuits, generally agreed to be the most "pure" form of asynchronous logic that retains computational universality. Less pure forms of asynchronous circuitry offer better performance at the cost of compromising one or more of these advantages:
Robust handling of metastability of arbiters. Early Completion of a circuit when it is known that the inputs which have not yet arrived are irrelevant. Possibly lower power consumption because no transistor ever transitions unless it is performing useful computation (clock gating in synchronous designs is an imperfect approximation of this ideal). Also, clock drivers can be removed which can significantly reduce power consumption. However, when using certain encodings, asynchronous circuits may require more area, which can result in increased power consumption if the underlying process has poor leakage properties (for example, deep submicrometer processes used prior to the introduction of high-K dielectrics). Freedom from the ever-worsening difficulties of distributing a high-fanout, timing-sensitive clock signal. Better modularity and composability. Far fewer assumptions about the manufacturing process are required (most assumptions are timing assumptions). Circuit speed is adapted on the fly to changing temperature and voltage conditions rather than being locked at the speed mandated by worst-case assumptions. Immunity to transistor-to-transistor variability in the manufacturing process, which is one of the most serious problems facing the semiconductor industry as dies shrink. Less severe electromagnetic interference. Synchronous circuits create a great deal of EMI in the frequency band at (or very near) their clock frequency and its harmonics; asynchronous circuits generate EMI patterns which are much more evenly spread across the spectrum. In asynchronous circuits, local signaling eliminates the need for global synchronization which exploits some potential advantages in comparison with synchronous ones. They have shown potential specifications in low power consumption, design reuse, improved noise immunity and electromagnetic compatibility. Asynchronous circuits are more tolerant to process variations and external voltage fluctuations‎[1]. Less stress on the power distribution network. Synchronous circuits tend to draw a large amount of current right at the clock edge and shortly thereafter. The number of nodes switching (and thence, amount of current drawn) drops off rapidly after the clock edge, reaching zero just before the next clock edge. In an asynchronous circuit, the switching times of the nodes are not correlated in this manner, so the current draw tends to be more uniform and less bursty.
Disadvantages:-
Requires people experienced in synchronous design to learn a new style. Performance analysis of asynchronous circuits may be challenging.

Application:-
Asynchronous CPUs are one of several ideas for radically changing CPU design.
Unlike a conventional processor, a clockless processor (asynchronous CPU) has no central clock to coordinate the progress of data through the pipeline. Instead, stages of the CPU are coordinated using logic devices called "pipeline controls" or "FIFO sequencers." Basically, the pipeline controller clocks the next stage of logic when the existing stage is complete. In this way, a central clock is unnecessary. It may actually be even easier to implement high performance devices in asynchronous, as opposed to clocked, logic:
components can run at different speeds on an asynchronous CPU; all major components of a clocked CPU must remain synchronized with the central clock; a traditional CPU cannot "go faster" than the expected worst-case performance of the slowest stage/instruction/component. When an asynchronous CPU completes an operation more quickly than anticipated, the next stage can immediately begin processing the results, rather than waiting for synchronization with a central clock. An operation might finish faster than normal because of attributes of the data being processed (e.g., multiplication can be very fast when multiplying by 0 or 1, even when running code produced by a naive compiler), or because of the presence of a higher voltage or bus speed setting, or a lower ambient temperature, than 'normal' or expected. Asynchronous logic proponents believe these capabilities would have these benefits:
lower power dissipation for a given performance level, and highest possible execution speeds. The biggest disadvantage of the clockless CPU is that most CPU design tools assume a clocked CPU (i.e., a synchronous circuit). Many tools "enforce synchronous design practices"[1]. Making a clockless CPU (designing an asynchronous circuit) involves modifying the design tools to handle clockless logic and doing extra testing to ensure the design avoids metastable problems. The group that designed the AMULET, for example, developed a tool called LARD to cope with the complex design of AMULET3.
Despite the difficulty of doing so, numerous asynchronous CPUs have been built, including:
the ORDVAC (?) and the (identical) ILLIAC I (1951), [2] the ILLIAC II (1962); The Caltech Asynchronous Microprocessor, the world-first asynchronous microprocessor (1988); the ARM-implementing AMULET (1993 and 2000); the asynchronous implementation of MIPS R3000, dubbed MiniMIPS (1998); an ARM-compatible processor (2003?) designed by Z. C. Yu, S. B. Furber, and L. A. Plana; "designed specifically to explore the benefits of asynchronous design for security sensitive applications";[3] the "Network-based Asynchronous Architecture" processor (2005) that executes a subset of the MIPS architecture instruction set;[3] the SEAforth multi-core processor (2008) from Charles H. Moore.[4] The ILLIAC II was the first completely asynchronous, speed independent processor design ever built; it was the most powerful computing machine known to man at the time.
DEC PDP-16 Register Transfer Modules (ca. 1973) allowed the experimenter to construct asynchronous, 16-bit processing elements. Delays for each module were fixed and based on the module's worst-case timing.
The Caltech Asynchronous Microprocessor (1988) was the first asynchronous microprocessor (1988). Caltech designed and manufactured the world's first fully Quasi Delay Insensitive processor.[citation needed] During demonstrations, the researchers amazed viewers by loading a simple program which ran in a tight loop, pulsing one of the output lines after each instruction. This output line was connected to an oscilloscope. When a cup of hot coffee was placed on the chip, the pulse rate (the effective "clock rate") naturally slowed down to adapt to the worsening performance of the heated transistors. When liquid nitrogen was poured on the chip, the instruction rate shot up with no additional intervention. Additionally, at lower temperatures, the voltage supplied to the chip could be safely increased, which also improved the instruction rate—again, with no additional configuration.
In 2004, Epson manufactured the world's first flexible microprocessor called ACT11, an 8-bit asynchronous chip. Synchronous flexible processors are slower, since bending the material on which a chip is fabricated causes wild and unpredictable variations in the delays of various transistors, for which worst case scenarios must be assumed everywhere and everything must be clocked at worst case speed. The processor is intended for use in smart cards, whose chips are currently limited in size to those small enough that they can remain perfectly rigid.

Ammeter

2009-11-08T09:35:00.000-08:00

An ammeter is a measuring instrument used to measure the electric current in a circuit. Electric currents are measured in amperes (A), hence the name. Smaller values of current can be measured using a milliameter or a microammeter. Early ammeters were laboratory instruments only which relied on the Earth's magnetic field for operation. By the late 19th century, improved instruments were designed which could be mounted in any position and allowed accurate measurements in electric power systems.
History:-
The relation between electric currents, magnetic fields and physical forces was first noted by Hans Christian Ørsted who in 1820 observed a compass needle was deflected from pointing North when a current flowed in an adjacent wire. The tangent galvanometer was used to measure currents using this effect, where the restoring force returning the pointer to the zero position was provided by the Earth's magnetic field. This made these instruments usable only when aligned with the Earth's field. Sensitivity of the instrument was increased by using additional turns of wire to multiply the effect - the instruments were called "multipliers".
Type:-

The D'Arsonval galvanometer is a moving coil ammeter. It uses magnetic deflection, where current passing through a coil causes the coil to move in a magnetic field. The voltage drop across the coil is kept to a minimum to minimize resistance across the ammeter in any circuit into which it is inserted. The modern form of this instrument was developed by Edward Weston, and uses two spiral springs to provide the restoring force. By maintaining a uniform air gap between the iron core of the instrument and the poles of its permanent magnet, the instrument has good linearity and accuracy. Basic meter movements can have full-scale deflection for currents from about 25 microamperes to 10 millamperes and have linear scales.
Moving iron ammeters use a piece of iron which moves when acted upon by the electromagnetic force of a fixed coil of wire. This type of meter responds to both direct and alternating currents (as opposed to the moving coil ammeter, which works on direct current only). The iron element consists of a moving vane attached to a pointer, and a fixed vane, surrounded by a coil. As alternating or direct current flows through the coil and induces a magnetic field in both vanes. The vanes repel each other and the moving vane deflects against the restoring force provided by fine helical springs.
An electrodynamic movement uses an electromagnet instead of the permanent magnet of the d'Arsonval movement. This instrument can respond to both alternating and direct current.
In a hot-wire ammeter, a current passes through a wire which expands as it heats. Although these instruments have slow response time and low accuracy, they are sometimes useful in measuring radio-frequency current.
Digital ammeter designs use an analog to digital converter (ADC) to measure the voltage across the shunt resistor; the digital display is calibrated to read the current through the shunt.
Application:-
To measure larger currents, a resistor called a shunt is placed in parallel with the meter. Most of the current flows through the shunt, and only a small fraction flows through the meter. This allows the meter to measure large currents. Traditionally, the meter used with a shunt has a full-scale deflection (FSD) of 50 mV, so shunts are typically designed to produce a voltage drop of 50 mV when carrying their full rated current.
Zero-center ammeters are used for applications requiring current to be measured with both polarities, common in scientific and industrial equipment. Zero-center ammeters are also commonly placed in series with a battery. In this application, the charging of the battery deflects the needle to one side of the scale (commonly, the right side) and the discharging of the battery deflects the needle to the other side.
Since the ammeter shunt has a very low resistance, mistakenly wiring the ammeter in parallel with a voltage source will cause a short circuit, at best blowing a fuse, possibly damaging the instrument and wiring, and exposing an observer to injury.
In AC circuits, a current transformer converts the magnetic field around a conductor into a small AC current, typically either 1 or 5 Amps at full rated current, that can be easily read by a meter. In a similar way, accurate AC/DC non-contact ammeters have been constructed using Hall effect magnetic field sensors. A portable hand-held clamp-on ammeter is a common tool for maintenance of industrial and commercial electrical equipment, which is temporarily clipped over a wire to measure current.

Electric charge

2009-11-07T09:46:00.000-08:00

Electric charge is a fundamental conserved property of some subatomic particles, which determines their electromagnetic interaction. Electrically charged matter is influenced by, and produces, electromagnetic fields. The interaction between a moving charge and an electromagnetic field is the source of the electromagnetic force, which is one of the four fundamental forces.
The electric charge on a body may be positive or negative. Two positively charged bodies experience a mutual repulsive force, as do two negatively charged bodies. A positively charged body and a negatively charged body experience an attractive force. The study of how charged bodies interact is classical electrodynamics, which is accurate insofar as quantum effects can be ignored.
Twentieth-century experiments demonstrated that electric charge is quantized: the charge of any system, body, or particle (except quarks) is an integer multiple of the elementary charge, e, approximately equal to 1.602×10−19 coulombs. The proton has a charge of e, and the electron has a charge of −e. The study of charged particles, and how their interactions are mediated by photons, is quantum electrodynamics

History:-
Electric charge is a characteristic property of many subatomic particles. The charges of free-standing particles are integer multiples of the elementary charge e; we say that electric charge is quantized. Michael Faraday, in his electrolysis experiments, was the first to note the discrete nature of electric charge. Robert Millikan's oil-drop experiment demonstrated this fact directly, and measured the elementary charge.
By convention, the charge of an electron is −1, while that of a proton is +1. Charged particles whose charges have the same sign repel one another, and particles whose charges have different signs attract. Coulomb's law quantifies the electrostatic force between two particles by asserting that the force is proportional to the product of their charges, and inversely proportional to the square of the distance between them.
The charge of an antiparticle equals that of the corresponding particle, but with opposite sign. Quarks have fractional charges of either −1⁄3 or +2⁄3, but free-standing quarks have never been observed (the theoretical reason for this fact is asymptotic freedom).
The electric charge of a macroscopic object is the sum of the electric charges of the particles that make it up. This charge is often zero, because matter is made of atoms, and atoms all have equal numbers of protons and electrons. More generally, in every molecule, the number of anions (negatively charged atoms) equals the number of cations (positively charged atoms). When the net electric charge is non-zero and motionless, the phenomenon is known as static electricity. Even when the net charge is zero, it can be distributed non-uniformly (e.g., due to an external electric field or to molecular motion), in which case the material is said to be polarized. The charge due to polarization is known as bound charge, while the excess charge brought from outside is called free charge. The motion of charged particles (especially the motion of electrons in metals) in a given direction is known as electric current.The SI unit of quantity of electric charge is the coulomb, which is equivalent to about 6.25 × 1018 e (e is the charge on a single electron or proton). Hence, the charge of an electron is approximately −1.602×10−19 C. The coulomb is defined as the quantity of charge that has passed through the cross-section of an electrical conductor carrying one ampere within one second. The symbol Q is often used to denote a quantity of electricity or charge. The quantity of electric charge can be directly measured with an electrometer, or indirectly measured with a ballistic galvanometer.
After finding the quantized character of charge, in 1891 Stoney proposed the unit 'electron' for this fundamental unit of electrical charge. This was before the discovery of the particle by J.J. Thomson in 1897. Today, the name "electron" for the unit of charge is no longer widely used except in the derived unit "electronvolt". This is quite surprising considering the wide use of this unit in the fields of physics and chemistry. The unit is today treated as nameless, referred to as "fundamental unit of charge" or simply as "e".
Formally, a measure of charge should be a multiple of the elementary charge e (charge is quantized), but since it is an average, macroscopic quantity, many orders of magnitude larger than a single elementary charge, it can effectively take on any real value. Furthermore, in some contexts it is meaningful to speak of fractions of a charge; e.g. in the charging of a capacitor.
As reported by the Ancient Greek philosopher Thales of Miletus around 600 BC, charge (or electricity) could be accumulated by rubbing fur on various substances, such as amber. The Greeks noted that the charged amber buttons could attract light objects such as hair. They also noted that if they rubbed the amber for long enough, they could even get a spark to jump. This property derives from the triboelectric effect.
In 1600 the English scientist William Gilbert returned to the subject in De Magnete, and coined the New Latin word electricus from ηλεκτρον (elektron), the Greek word for "amber", which soon gave rise to the English words "electric" and "electricity." He was followed in 1660 by Otto von Guericke, who invented what was probably the first electrostatic generator. Other European pioneers were Robert Boyle, who in 1675 stated that electric attraction and repulsion can act across a vacuum; Stephen Gray, who in 1729 classified materials as conductors and insulators; and C. F. du Fay, who proposed in 1733[1] that electricity came in two varieties which cancelled each other, and expressed this in terms of a two-fluid theory. When glass was rubbed with silk, du Fay said that the glass was charged with vitreous electricity, and when amber was rubbed with fur, the amber was said to be charged with resinous electricity. In 1839, Michael Faraday showed that the apparent division between static electricity, current electricity and bioelectricity was incorrect, and all were a consequence of the behavior of a single kind of electricity appearing in opposite polarities. It is arbitrary which polarity you call positive and which you call negative. Positive charge can be defined as the charge left on a glass rod after being rubbed with silk.
One of the foremost experts on electricity in the 18th century was Benjamin Franklin, who argued in favour of a one-fluid theory of electricity. Franklin imagined electricity as being a type of invisible fluid present in all matter; for example he believed that it was the glass in a Leyden jar that held the accumulated charge. He posited that rubbing insulating surfaces together caused this fluid to change location, and that a flow of this fluid constitutes an electric current. He also posited that when matter contained too little of the fluid it was "negatively" charged, and when it had an excess it was "positively" charged. Arbitrarily (or for a reason that was not recorded) he identified the term "positive" with vitreous electricity and "negative" with resinous electricity. William Watson arrived at the same explanation at about the same time.

Digital circuit

2009-11-05T06:26:00.000-08:00

Digital electronics are systems that represent signals as discrete levels, rather than as a continuous range. In most cases the number of states is two, and these states are represented by two voltage levels: one near to zero volts and one at a higher level depending on the supply voltage in use. These two levels are often represented as "Low" and "High."
The fundamental advantage of digital techniques stem from the fact it is easier to get an electronic device to switch into one of a number of known states than to accurately reproduce a continuous range of values.Digital electronics are usually made from large assemblies of logic gates, simple electronic representations of Boolean logic functions
Advantages:-
One advantage of digital circuits when compared to analog circuits is that signals represented digitally can be transmitted without degradation due to noise. For example, a continuous audio signal, transmitted as a sequence of 1s and 0s, can be reconstructed without error provided the noise picked up in transmission is not enough to prevent identification of the 1s and 0s. An hour of music can be stored on a compact disc as about 6 billion binary digits.
In a digital system, a more precise representation of a signal can be obtained by using more binary digits to represent it. While this requires more digital circuits to process the signals, each digit is handled by the same kind of hardware. In an analog system, additional resolution requires fundamental improvements in the linearity and noise charactersitics of each step of the signal chain. Computer-controlled digital systems can be controlled by software, allowing new functions to be added without changing hardware. Often this can be done outside of the factory by updating the product's software. So, the product's design errors can be corrected after the product is in a customer's hands. Information storage can be easier in digital systems than in analog ones. The noise-immunity of digital systems permits data to be stored and retrieved without degradation. In an analog system, noise from aging and wear degrade the information stored. In a digital system, as long as the total noise is below a certain level, the information can be recovered perfectly.
Disadvantages:-
In some cases, digital circuits use more energy than analog circuits to accomplish the same tasks, thus producing more heat. In portable or battery-powered systems this can limit use of digital systems. For example, battery-powered cellular telephones often use a low-power analog front-end to amplify and tune in the radio signals from the base station. However, a base station has grid power and can use power-hungry, but very flexible software radios. Such base stations can be easily reprogrammed to process the signals used in new cellular standards.
Digital circuits are sometimes more expensive, especially in small quantities.
The sensed world is analog, and signals from this world are analog quantities. For example, light, temperature, sound, electrical conductivity, electric and magnetic fields are analog. Most useful digital systems must translate from continuous analog signals to discrete digital signals. This causes quantization errors. Quantization error can be reduced if the system stores enough digital data to represent the signal to the desired degree of fidelity. The Nyquist-Shannon sampling theorem provides an important guideline as to how much digital data is needed to accurately portray a given analog signal. In some systems, if a single piece of digital data is lost or misinterpreted, the meaning of large blocks of related data can completely change. Because of the cliff effect, it can be difficult for users to tell if a particular system is right on the edge of failure, or if it can tolerate much more noise before failing.
Digital fragility can be reduced by designing a digital system for robustness. For example, a parity bit or other error management method can be inserted into the signal path. These schemes help the system detect errors, and then either correct the errors, or at least ask for a new copy of the data. In a state-machine, the state transition logic can be designed to catch unused states and trigger a reset sequence or other error recovery routine.
Embedded software designs that employ Immunity Aware Programming, such as the practice of filling unused program memory with interrupt instructions that point to an error recovery routine. This helps guard against failures that corrupt the microcontroller's instruction pointer which could otherwise cause random code to be executed. Digital memory and transmission systems can use techniques such as error detection and correction to use additional data to correct any errors in transmission and storage. On the other hand, some techniques used in digital systems make those systems more vulnerable to single-bit errors. These techniques are acceptable when the underlying bits are reliable enough that such errors are highly unlikely.
A single-bit error in audio data stored directly as linear pulse code modulation (such as on a CD-ROM) causes, at worst, a single click. Instead, many people use audio compression to save storage space and download time, even though a single-bit error may corrupt the entire song.
Structure of digital system:-Engineers use many methods to minimize logic functions, in order to reduce the circuit's complexity. When the complexity is less, the circuit also has fewer errors and less electronics, and is therefore less expensive.
The most widely used simplification is a minimization algorithm like the Espresso heuristic logic minimizer within a CAD system, although historically, binary decision diagrams, an automated Quine–McCluskey algorithm, truth tables, Karnaugh Maps, and Boolean algebra have been used. Representations are crucial to an engineer's design of digital circuits. Some analysis methods only work with particular representations.
The classical way to represent a digital circuit is with an equivalent set of logic gates. Another way, often with the least electronics, is to construct an equivalent system of electronic switches (usually transistors). One of the easiest ways is to simply have a memory containing a truth table. The inputs are fed into the address of the memory, and the data outputs of the memory become the outputs.
For automated analysis, these representations have digital file formats that can be processed by computer programs. Most digital engineers are very careful to select computer programs ("tools") with compatible file formats.
To choose representations, engineers consider types of digital systems. Most digital systems divide into "combinational systems" and "sequential systems." A combinational system always presents the same output when given the same inputs. It is basically a representation of a set of logic functions, as already discussed.
A sequential system is a combinational system with some of the outputs fed back as inputs. This makes the digital machine perform a "sequence" of operations. The simplest sequential system is probably a flip flop, a mechanism that represents a binary digit or "bit".
Sequential systems are often designed as state machines. In this way, engineers can design a system's gross behavior, and even test it in a simulation, without considering all the details of the logic functions.
Sequential systems divide into two further subcategories. "Synchronous" sequential systems change state all at once, when a "clock" signal changes state. "Asynchronous" sequential systems propagate changes whenever inputs change. Synchronous sequential systems are made of well-characterized asynchronous circuits such as flip-flops, that change only when the clock changes, and which have carefully designed timing margins.
The usual way to implement a synchronous sequential state machine is divide it into a piece of combinational logic and a set of flip flops called a "state register." Each time a clock signal ticks, the state register captures the feedback generated from the previous state of the combinational logic, and feeds it back as an unchanging input to the combinational part of the state machine. The fastest rate of the clock is set by the most time-consuming logic calculation in the combinational logic.
The state register is just a representation of a binary number. If the states in the state machine are numbered (easy to arrange), the logic function is some combinational logic that produces the number of the next state.
In comparison, asynchronous systems are very hard to design because all possible states, in all possible timings must be considered. The usual method is to construct a table of the minimum and maximum time that each such state can exist, and then adjust the circuit to minimize the number of such states, and force the circuit to periodically wait for all of its parts to enter a compatible state. (This is called "self-resynchronization.") Without such careful design, it is easy to accidentally produce asynchronous logic that is "unstable", that is, real electronics will have unpredictable results because of the cumulative delays caused by small variations in the values of the electronic components. Certain circuits (such as the synchronizer flip-flops, switch debouncers, and the like which allow external unsynchronized signals to enter synchronous logic circuits) are inherently asynchronous in their design and must be analyzed as such.
As of 2005, almost all digital machines are synchronous designs because it is much easier to create and verify a synchronous design—the software currently used to simulate digital machines does not yet handle asynchronous designs. However, asynchronous logic is thought to be superior, if it can be made to work, because its speed is not constrained by an arbitrary clock; instead, it simply runs at the maximum speed permitted by the propagation rates of the logic gates from which it is constructed. Building an asynchronous circuit using faster parts implicitly makes the circuit "go" faster.
More generally, many digital systems are data flow machines. These are usually designed using synchronous register transfer logic, using hardware description languages such as VHDL or Verilog.
In register transfer logic, binary numbers are stored in groups of flip flops called registers. The outputs of each register are a bundle of wires called a "bus" that carries that number to other calculations. A calculation is simply a piece of combinational logic. Each calculation also has an output bus, and these may be connected to the inputs of several registers. Sometimes a register will have a multiplexer on its input, so that it can store a number from any one of several buses. Alternatively, the outputs of several items may be connected to a bus through buffers that can turn off the output of all of the devices except one. A sequential state machine controls when each register accepts new data from its input.
In the 1980s, some researchers discovered that almost all synchronous register-transfer machines could be converted to asynchronous designs by using first-in-first-out synchronization logic. In this scheme, the digital machine is characterized as a set of data flows. In each step of the flow, an asynchronous "synchronization circuit" determines when the outputs of that step are valid, and presents a signal that says, "grab the data" to the stages that use that stage's inputs. It turns out that just a few relatively simple synchronization circuits are needed.
The most general-purpose register-transfer logic machine is a computer. This is basically an automatic binary abacus. The control unit of a computer is usually designed as a microprogram run by a microsequencer. A microprogram is much like a player-piano roll. Each table entry or "word" of the microprogram commands the state of every bit that controls the computer. The sequencer then counts, and the count addresses the memory or combinational logic machine that contains the microprogram. The bits from the microprogram control the arithmetic logic unit, memory and other parts of the computer, including the microsequencer itself.
In this way, the complex task of designing the controls of a computer is reduced to a simpler task of programming a relatively independent collection of much simpler logic machines.
Computer architecture is a specialized engineering activity that tries to arrange the registers, calculation logic, buses and other parts of the computer in the best way for some purpose. Computer architects have applied large amounts of ingenuity to computer design to reduce the cost and increase the speed and immunity to programming errors of computers. An increasingly common goal is to reduce the power used in a battery-powered computer system, such as a cell-phone. Many computer architects serve an extended apprenticeship as microprogrammers.
"Specialized computers" are usually a conventional computer with a special-purpose microprogram.
Automated design tools:-
To save costly engineering effort, much of the effort of designing large logic machines has been automated. The computer programs are called "electronic design automation tools" or just "EDA."
Simple truth table-style descriptions of logic are often optimized with EDA that automatically produces reduced systems of logic gates or smaller lookup tables that still produce the desired outputs. The most common example of this kind of software is the Espresso heuristic logic minimizer.Most practical algorithms for optimizing large logic systems use algebraic manipulations or binary decision diagrams, and there are promising experiments with genetic algorithms and annealing optimizations.To automate costly engineering processes, some EDA can take state tables that describe state machines and automatically produce a truth table or a function table for the combinatorial part of a state machine. The state table is a piece of text that lists each state, together with the conditions controlling the transitions between them and the belonging output signals.
It is common for the function tables of such computer-generated state-machines to be optimized with logic-minimization software such as Minilog.
Often, real logic systems are designed as a series of sub-projects, which are combined using a "tool flow." The tool flow is usually a "script," a simplified computer language that can invoke the software design tools in the right order.Tool flows for large logic systems such as microprocessors can be thousands of commands long, and combine the work of hundreds of engineers.Writing and debugging tool flows is an established engineering specialty in companies that produce digital designs. The tool flow usually terminates in a detailed computer file or set of files that describe how to physically construct the logic. Often it consists of instructions to draw the transistors and wires on an integrated circuit or a printed circuit board.
Parts of tool flows are "debugged" by verifying the outputs of simulated logic against expected inputs. The test tools take computer files with sets of inputs and outputs, and highlight discrepancies between the simulated behavior and the expected behavior.
Once the input data is believed correct, the design itself must still be verified for correctness. Some tool flows verify designs by first producing a design, and then scanning the design to produce compatible input data for the tool flow. If the scanned data matches the input data, then the tool flow has probably not introduced errors.
The functional verification data are usually called "test vectors." The functional test vectors may be preserved and used in the factory to test that newly constructed logic works correctly. However, functional test patterns don't discover common fabrication faults. Production tests are often designed by software tools called "test pattern generators." These generate test vectors by examining the structure of the logic and systematically generating tests for particular faults. This way the fault coverage can closely approach 100%, provided the design is properly made testable (see next section).Once a design exists, and is verified and testable, it often needs to be processed to be manufacturable as well. Modern integrated circuits have features smaller than the wavelength of the light used to expose the photoresist. Manufacturability software adds interference patterns to the exposure masks to eliminate open-circuits, and enhance the masks' resolution and contrast.

Magnetism

2009-11-03T07:48:00.000-08:00

This article is about magnetic materials. For information about objects and devices that produce a magnetic field, see magnet. For field that magnets and currents produce, see magnetic field. For other uses, see magnetism (disambiguation).In physics, the term magnetism is used to describe how materials respond on the microscopic level to an applied magnetic field; to categorize the magnetic phase of a material. For example, the most well known form of magnetism is ferromagnetism such that some ferromagnetic materials produce their own persistent magnetic field. However, all materials are influenced to greater or lesser degree by the presence of a magnetic field. Some are attracted to a magnetic field (paramagnetism); others are repulsed by a magnetic field (diamagnetism); others have a much more complex relationship with an applied magnetic field. Substances that are negligibly affected by magnetic fields are known as non-magnetic substances. They include copper, aluminium, water, and gases.
The magnetic state (or phase) of a material depends on temperature (and other variables such as pressure and applied magnetic field) so that a material may exhibit more than one form of magnetism depending on its temperature, etc.

History:-
Aristotle attributes the first of what could be called a scientific discussion on magnetism to Thales, who lived from about 625 BC to about 545 BC Around the same time in ancient India, the Indian surgeon, Sushruta, was the first to make use of the magnet for surgical purposes
In ancient China, the earliest literary reference to magnetism lies in a 4th century BC book called Book of the Devil Valley Master (鬼谷子): "The lodestone makes iron come or it attracts it."The earliest mention of the attraction of a needle appears in a work composed between AD 20 and 100 (Louen-heng): "A lodestone attracts a needle." The ancient Chinese scientist Shen Kuo (1031-1095) was the first person to write of the magnetic needle compass and that it improved the accuracy of navigation by employing the astronomical concept of true north (Dream Pool Essays, AD 1088 ), and by the 12th century the Chinese were known to use the lodestone compass for navigation.
Alexander Neckham, by 1187, was the first in Europe to describe the compass and its use for navigation. In 1269, Peter Peregrinus de Maricourt wrote the Epistola de magnete, the first extant treatise describing the properties of magnets. In 1282, the properties of magnets and the dry compass were discussed by Al-Ashraf, a Yemeni physicist, astronomer and geographer.
In 1600, William Gilbert published his De Magnete, Magneticisque Corporibus, et de Magno Magnete Tellure (On the Magnet and Magnetic Bodies, and on the Great Magnet the Earth). In this work he describes many of his experiments with his model earth called the terrella. From his experiments, he concluded that the Earth was itself magnetic and that this was the reason compasses pointed north (previously, some believed that it was the pole star (Polaris) or a large magnetic island on the north pole that attracted the compass).
An understanding of the relationship between electricity and magnetism began in 1819 with work by Hans Christian Oersted, a professor at the University of Copenhagen, who discovered more or less by accident that an electric current could influence a compass needle. This landmark experiment is known as Oersted's Experiment. Several other experiments followed, with André-Marie Ampère, Carl Friedrich Gauss, Michael Faraday, and others finding further links between magnetism and electricity. James Clerk Maxwell synthesized and expanded these insights into Maxwell's equations, unifying electricity, magnetism, and optics into the field of electromagnetism. In 1905, Einstein used these laws in motivating his theory of special relativity, requiring that the laws held true in all inertial reference frames.
Electromagnetism has continued to develop into the twenty-first century, being incorporated into the more fundamental theories of gauge theory, quantum electrodynamics, electroweak theory, and finally the standard model.

Sources of magnetism:-
There exists a close connection between angular momentum and magnetism, expressed on a macroscopic scale in the Einstein-de Haas effect "rotation by magnetization" and its inverse, the Barnett effect or "magnetization by rotation".At the atomic and sub-atomic scales, this connection is expressed by the ratio of magnetic moment to angular momentum, the gyromagnetic ratio.Magnetism, at its root, arises from two sources:Electric currents, or more generally moving electric charges, create magnetic fields (see Maxwell's Equations). Many particles have nonzero "intrinsic" (or "spin") magnetic moments. (Just as each particle, by its nature, has a certain mass and charge, each has a certain magnetic moment, possibly zero.) In magnetic materials, the most important sources of magnetization are, more specifically, the electrons' orbital angular motion around the nucleus, and the electrons' intrinsic magnetic moment (see Electron magnetic dipole moment). The other potential sources of magnetism are much less important: For example, the nuclear magnetic moments of the nuclei in the material are typically thousands of times smaller than the electrons' magnetic moments, so they are negligible in the context of the magnetization of materials. (Nuclear magnetic moments are important in other contexts, particularly in Nuclear Magnetic Resonance (NMR) and Magnetic Resonance Imaging (MRI).)
Ordinarily, the countless electrons in a material are arranged such that their magnetic moments (both orbital and intrinsic) cancel out. This is due, to some extent, to electrons combining into pairs with opposite intrinsic magnetic moments (as a result of the Pauli exclusion principle; see Electron configuration), or combining into "filled subshells" with zero net orbital motion; in both cases, the electron arrangement is so as to exactly cancel the magnetic moments from each electron. Moreover, even when the electron configuration is such that there are unpaired electrons and/or non-filled subshells, it is often the case that the various electrons in the solid will contribute magnetic moments that point in different, random directions, so that the material will not be magnetic.However, sometimes (either spontaneously, or owing to an applied external magnetic field) each of the electron magnetic moments will be, on average, lined up. Then the material can produce a net total magnetic field, which can potentially be quite strong.
The magnetic behavior of a material depends on its structure (particularly its electron configuration, for the reasons mentioned above), and also on the temperature (at high temperatures, random thermal motion makes it more difficult for the electrons to maintain alignment).

Laser diode

2009-11-02T08:46:00.000-08:00

A laser diode is a laser where the active medium is a semiconductor similar to that found in a light-emitting diode. The most common and practical type of laser diode is formed from a p-n junction and powered by injected electric current. These devices are sometimes referred to as injection laser diodes to distinguish them from (optically) pumped laser diodes, which are more easily manufactured in the laboratory.
Theory of operation:-
A laser diode, like many other semiconductor devices, is formed by doping a very thin layer on the surface of a crystal wafer. The crystal is doped to produce an n-type region and a p-type region, one above the other, resulting in a p-n junction, or diode.
The many types of diode lasers known today collectively form a subset of the larger classification of semiconductor p-n junction diodes. Just as in any semiconductor p-n junction diode, forward electrical bias causes the two species of charge carrier - holes and electrons - to be "injected" from opposite sides of the p-n junction into the depletion region, situated at its heart. Holes are injected from the p-doped, and electrons from the n-doped, semiconductor. (A depletion region, devoid of any charge carriers, forms automatically and unavoidably as a result of the difference in chemical potential between n- and p-type semiconductors wherever they are in physical contact.)As charge injection is a distinguishing feature of diode lasers as compared to all other lasers, diode lasers are traditionally and more formally called "injection lasers." (This terminology differentiates diode lasers, e.g., from flashlamp-pumped solid state lasers, such as the ruby laser. Interestingly, whereas the term "solid-state" was extremely apt in differentiating 1950s-era semiconductor electronics from earlier generations of vacuum electronics, it would not have been adequate to convey unambiguously the unique characteristics defining 1960s-era semiconductor lasers.) When an electron and a hole are present in the same region, they may recombine or "annihilate" with the result being spontaneous emission — i.e., the electron may re-occupy the energy state of the hole, emitting a photon with energy equal to the difference between the electron and hole states involved. (In a conventional semiconductor junction diode, the energy released from the recombination of electrons and holes is carried away as phonons, i.e., lattice vibrations, rather than as photons.) Spontaneous emission gives the laser diode below lasing threshold similar properties to an LED. Spontaneous emission is necessary to initiate laser oscillation, but it is one among several sources of inefficiency once the laser is oscillating.
The difference between the photon-emitting semiconductor laser (or LED) and conventional phonon-emitting (non-light-emitting) semiconductor junction diodes lies in the use of a different type of semiconductor, one whose physical and atomic structure confers the possibility for photon emission. These photon-emitting semiconductors are the so-called "direct bandgap" semiconductors. The properties of silicon and germanium, which are single-element semiconductors, have bandgaps that do not align in the way needed to allow photon emission and are not considered "direct." Other materials, the so-called compound semiconductors, have virtually identical crystaline structures as silicon or germanium but use alternating arrangements of two different atomic species in a checkerboard-like pattern to break the symmetry. The transition between the materials in the alternating pattern creates the critical "direct bandgap" property. Gallium arsenide, indium phosphide, gallium antimonide, and gallium nitride are all examples of compound semiconductor materials that can be used to create junction diodes that emit light.
In the absence of stimulated emission (e.g., lasing) conditions, electrons and holes may coexist in proximity to one another, without recombining, for a certain time, termed the "upper-state lifetime" or "recombination time" (about a nanosecond for typical diode laser materials), before they recombine. Then a nearby photon with energy equal to the recombination energy can cause recombination by stimulated emission. This generates another photon of the same frequency, travelling in the same direction, with the same polarization and phase as the first photon. This means that stimulated emission causes gain in an optical wave (of the correct wavelength) in the injection region, and the gain increases as the number of electrons and holes injected across the junction increases. The spontaneous and stimulated emission processes are vastly more efficient in direct bandgap semiconductors than in indirect bandgap semiconductors; therefore silicon is not a common material for laser diodes.
As in other lasers, the gain region is surrounded with an optical cavity to form a laser. In the simplest form of laser diode, an optical waveguide is made on that crystal surface, such that the light is confined to a relatively narrow line. The two ends of the crystal are cleaved to form perfectly smooth, parallel edges, forming a Fabry-Perot resonator. Photons emitted into a mode of the waveguide will travel along the waveguide and be reflected several times from each end face before they are emitted. As a light wave passes through the cavity, it is amplified by stimulated emission, but light is also lost due to absorption and by incomplete reflection from the end facets. Finally, if there is more amplification than loss, the diode begins to "lase".
Some important properties of laser diodes are determined by the geometry of the optical cavity. Generally, in the vertical direction, the light is contained in a very thin layer, and the structure supports only a single optical mode in the direction perpendicular to the layers. In the lateral direction, if the waveguide is wide compared to the wavelength of light, then the waveguide can support multiple lateral optical modes, and the laser is known as "multi-mode". These laterally multi-mode lasers are adequate in cases where one needs a very large amount of power, but not a small diffraction-limited beam; for example in printing, activating chemicals, or pumping other types of lasers.
In applications where a small focused beam is needed, the waveguide must be made narrow, on the order of the optical wavelength. This way, only a single lateral mode is supported and one ends up with a diffraction-limited beam. Such single spatial mode devices are used for optical storage, laser pointers, and fiber optics. Note that these lasers may still support multiple longitudinal modes, and thus can lase at multiple wavelengths simultaneously.
The wavelength emitted is a function of the band-gap of the semiconductor and the modes of the optical cavity. In general, the maximum gain will occur for photons with energy slightly above the band-gap energy, and the modes nearest the gain peak will lase most strongly. If the diode is driven strongly enough, additional side modes may also lase. Some laser diodes, such as most visible lasers, operate at a single wavelength, but that wavelength is unstable and changes due to fluctuations in current or temperature.
Due to diffraction, the beam diverges (expands) rapidly after leaving the chip, typically at 30 degrees vertically by 10 degrees laterally. A lens must be used in order to form a collimated beam like that produced by a laser pointer. If a circular beam is required, cylindrical lenses and other optics are used. For single spatial mode lasers, using symmetrical lenses, the collimated beam ends up being elliptical in shape, due to the difference in the vertical and lateral divergences. This is easily observable with a red laser pointer.
The simple diode described above has been heavily modified in recent years to accommodate modern technology, resulting in a variety of types of laser diodes, as described below.
Laser diode types:-
The simple laser diode structure, described above, is extremely inefficient. Such devices require so much power that they can only achieve pulsed operation without damage. Although historically important and easy to explain, such devices are not practical.
Applications of laser diodes:-
Laser diodes are numerically the most common type of laser, with 2004 sales of approximately 733 million diode lasers, as compared to 131,000 of other types of lasers.Laser diodes find wide use in telecommunication as easily modulated and easily coupled light sources for fiber optics communication. They are used in various measuring instruments, eg. rangefinders. Another common use is in barcode readers. Visible lasers, typically red but later also green, are common as laser pointers. Both low and high-power diodes are used extensively in the printing industry both as light sources for scanning (input) of images and for very high-speed and high-resolution printing plate (output) manufacturing. Infrared and red laser diodes are common in CD players, CD-ROMs and DVD technology. Violet lasers are used in HD DVD and Blu-ray technology. Diode lasers have also found many applications in laser absorption spectrometry (LAS) for high-speed, low-cost assessment or monitoring of the concentration of various species in gas phase. High-power laser diodes are used in industrial applications such as heat treating, cladding, seam welding and for pumping other lasers, such as diode pumped solid state lasers.Applications of laser diodes can be categorized in various ways. Most applications could be served by larger solid state lasers or optical parametric oscillators, but the low cost of mass-produced diode lasers makes them essential for mass-market applications. Diode lasers can be used in a great many fields; since light has many different properties (power, wavelength and spectral quality, beam quality, polarization, etc.) it is interesting to classify applications by these basic properties.Many applications of diode lasers primarily make use of the "directed energy" property of an optical beam. In this category one might include the laser printers, bar-code readers, image scanning, illuminators, designators, optical data recording, combustion ignition, laser surgery, industrial sorting, industrial machining, and directed energy weaponry. Some of these applications are emerging while others are well-established.Laser medicine: medicine and especially dentistry have found many new applications for diode lasers. The shrinking size of the units and their increasing user friendliness makes them very attractive to clinicians for minor soft tissue procedures. The 800 nm - 980 nm units have a high absorption rate for hemoglobin and thus make them ideal for soft tissue applications, where good hemostasis is necessary.Applications which may today or in the future make use of the coherence of diode-laser-generated light include interferometric distance measurement, holography, coherent communications, and coherent control of chemical reactions.
Applications which may make use of "narrow spectral" properties of diode lasers include range-finding, telecommunications, infra-red countermeasures, spectroscopic sensing, generation of radio-frequency or terahertz waves, atomic clock state preparation, quantum key cryptography, frequency doubling and conversion, water purification (in the UV), and photodynamic therapy (where a particular wavelength of light would cause a substance such as porphyrin to become chemically active as an anti-cancer agent only where the tissue is illuminated by light).
Applications where the desired quality of laser diodes is their ability to generate ultra-short pulses of light by the technique known as "mode-locking" include clock distribution for high-performance integrated circuits, high-peak-power sources for laser-induced breakdown spectroscopy sensing, arbitrary waveform generation for radio-frequency waves, photonic sampling for analog-to-digital conversion, and optical code-division-multiple-access systems for secure communication.

Kirchhoff's circuit laws

2009-11-02T08:36:00.000-08:00

Kirchhoff's circuit laws are two equalities that deal with the conservation of charge and energy in electrical circuits, and were first described in 1845 by Gustav Kirchhoff. Widely used in electrical engineering, they are also called Kirchhoff's rules or simply Kirchhoff's laws (see also Kirchhoff's laws for other meanings of that term).
Both circuit rules can be directly derived from Maxwell's equations, but Kirchhoff preceded Maxwell and instead generalized work by Georg Ohm.

Kcl:-
This law is also called Kirchhoff's point rule, Kirchhoff's junction rule (or nodal rule), and Kirchhoff's first rule.The principle of conservation of electric charge implies that:
At any node (junction) in an electrical circuit, the sum of currents flowing into that node is equal to the sum of currents flowing out of that node. Adopting the convention that every current flowing towards the node is positive and that every current flowing away is negative (or the other way around), this principle can be stated as:
n is the total number of branches with currents flowing towards or away from the node.
This formula is also valid for complex currents:
The law is based on the conservation of charge whereby the charge (measured in coulombs) is the product of the current (in amps) and the time (which is measured in seconds).
Changing charge density:-
Physically speaking, the restriction regarding the "capacitor plate" means that Kirchhoff's current law is only valid if the charge density remains constant in the point that it is applied to. This is normally not a problem because of the strength of electrostatic forces: the charge buildup would cause repulsive forces to disperse the charges.
However, a charge build-up can occur in a capacitor, where the charge is typically spread over wide parallel plates, with a physical break in the circuit that prevents the positive and negative charge accumulations over the two plates from coming together and cancelling. In this case, the sum of the currents flowing into one plate of the capacitor is not zero, but rather is equal to the rate of charge accumulation. However, if the displacement current dD/dt is included, Kirchhoff's current law once again holds. (This is really only required if one wants to apply the current law to a point on a capacitor plate. In circuit analyses, however, the capacitor as a whole is typically treated as a unit, in which case the ordinary current law holds since exactly the current that enters the capacitor on the one side leaves it on the other side.)
More technically, Kirchhoff's current law can be found by taking the divergence of Ampère's law with Maxwell's correction and combining with Gauss's law, yielding:
This is simply the charge conservation equation (in integral form, it says that the current flowing out of a closed surface is equal to the rate of loss of charge within the enclosed volume (Divergence theorem)). Kirchhoff's current law is equivalent to the statement that the divergence of the current is zero, true for time-invariant ρ, or always true if the displacement current is included with J.

Coulomb's law

2009-11-02T08:32:00.000-08:00

Coulomb's law, sometimes called the Coulomb law, is an equation describing the electrostatic force between electric charges. It was studied and first published in the 1780s by French physicist Charles Augustin de Coulomb and was essential to the development of the theory of electromagnetism. Nevertheless, the dependence of the electric force with distance (inverse square law) had been proposed previously by Joseph Priestley and the dependence with both distance and charge had been discovered, but not published, by Henry Cavendish, prior to Coulomb's works.
Scalar form:-
The scalar form of Coulomb's law will only describe the magnitude of the electrostatic force between two electric charges. If direction is required, then the vector form is required as well. The magnitude of the electrostatic force (F) on a charge (q1) due to the presence of a second charge (q2), is given by
where r is the distance between the two charges and ke a proportionality constant. A positive force implies a repulsive interaction, while a negative force implies an attractive interaction.
The proportionality constant ke, called Coulomb's constant (sometimes called Coulomb's force constant) is related to the properties of space and can be calculated exactly:
In SI units the speed of light in vacuum, denoted c0[4] is defined as 299,792,458 m·s−1, and the magnetic constant (μ0), is defined as 4π × 10−7 H·m−1,[6] leading to the definition for the electric constant (ε0) as ε0 = 1/(μ0c20) ≈ 8.854187817×10−12 F·m−1. In cgs units, the unit charge, esu of charge or statcoulomb, is defined so that this Coulomb force constant is 1.
This formula says that the magnitude of the force is directly proportional to the magnitude of the charges of each object and inversely proportional to the square of the distance between them. The exponent in Coulomb's Law has been found to differ from −2 by less than one in a billion.
When measured in units that people commonly use (such as SI—see International System of Units), the electrostatic force constant (ke) is numerically much much larger than the universal gravitational constant (G).This means that for objects with charge that is of the order of a unit charge (C) and mass of the order of a unit mass (kg), the electrostatic forces will be so much larger than the gravitational forces that the latter force can be ignored. This is not the case when Planck units are used and both charge and mass are of the order of the unit charge and unit mass. However, charged elementary particles have mass that is far less than the Planck mass while their charge is about the Planck charge so that, again, gravitational forces can be ignored. For example, the electrostatic force between an electron and a proton, which constitute a hydrogen atom, is almost 40 orders of magnitude greater than the gravitational force between them.
Coulomb's law can also be interpreted in terms of atomic units with the force expressed in Hartrees per Bohr radius, the charge in terms of the elementary charge, and the distances in terms of the Bohr radius.

Eddy current

2009-11-01T10:03:00.000-08:00

An eddy current (also known as Foucault current) is an electrical phenomenon discovered by French physicist François Arago in 1824. It is caused when a conductor is exposed to a changing magnetic field due to relative motion of the field source and conductor; or due to variations of the field with time. This can cause a circulating flow of electrons, or a current, within the body of the conductor. These circulating eddies of current create induced magnetic fields that oppose the change of the original magnetic field due to Lenz's law, causing repulsive or drag forces between the conductor and the magnet. The stronger the applied magnetic field, or the greater the electrical conductivity of the conductor,or the faster the field that the conductor is exposed to changes, then the greater the currents that are developed and the greater the opposing field.
The term eddy current comes from analogous currents seen in water when dragging an oar breadthwise: localised areas of turbulence known as eddies give rise to persistent vortices.
eddy currents, like all electric currents, generate heat as well as electromagnetic forces. The heat can be harnessed for induction heating. The electromagnetic forces can be used for levitation, creating movement, or to give a strong braking effect. Eddy currents can often be minimised with thin plates, by lamination of conductors or other details of conductor shape.
Explanation:-
When a conductor moves relative to the field generated by a source, electromotive forces (EMFs) can be generated around loops within the conductor. These EMFs acting on the resistivity of the material generate a current around the loop, in accordance with Faraday's law of induction. These currents dissipate energy, and create a magnetic field that tends to oppose the changes in the field.
Eddy currents are created when a moving conductor experiences changes in the magnetic field generated by a stationary object, as well as when a stationary conductor encounters a varying magnetic field. Both effects are present when a conductor moves through a varying magnetic field, as is the case at the top and bottom edges of the magnetized region shown in the diagram. Eddy currents will be generated wherever a conducting object experiences a change in the intensity or direction of the magnetic field at any point within it, and not just at the boundaries.
The swirling current set up in the conductor is due to electrons experiencing a Lorentz force that is perpendicular to their motion. Hence, they veer to their right, or left, depending on the direction of the applied field and whether the strength of the field is increasing or declining. The resistivity of the conductor acts to damp the amplitude of the eddy currents, as well as straighten their paths. Lenz's law encapsulates the fact that the current swirls in such a way as to create an induced magnetic field that opposes the phenomenon that created it. In the case of a varying applied field, the induced field will always be in the opposite direction to that applied. The same will be true when a varying external field is increasing in strength. However, when a varying field is falling in strength, the induced field will be in the same direction as that originally applied, in order to oppose the decline.
An object or part of an object experiences steady field intensity and direction where there is still relative motion of the field and the object (for example in the center of the field in the diagram), or unsteady fields where the currents cannot circulate due to the geometry of the conductor. In these situations charges collect on or within the object and these charges then produce static electric potentials that oppose any further current. Currents may be initially associated with the creation of static potentials, but these may be transitory and small.
Eddy currents generate resistive losses that transform some forms of energy, such as kinetic energy, into heat. In many devices, this Joule heating reduces efficiency of iron-core transformers and electric motors and other devices that use changing magnetic fields. Eddy currents are minimized in these devices by selecting magnetic core materials that have low electrical conductivity (e.g., ferrites) or by using thin sheets of magnetic material, known as laminations. Electrons cannot cross the insulating gap between the laminations and so are unable to circulate on wide arcs. Charges gather at the lamination boundaries, in a process analogous to the Hall effect, producing electric fields that oppose any further accumulation of charge and hence suppressing the eddy currents. The shorter the distance between adjacent laminations (i.e., the greater the number of laminations per unit area, perpendicular to the applied field), the greater the suppression of eddy currents.
The conversion of input energy to heat is not always undesirable, however, as there are some practical applications. One is in the brakes of some trains known as eddy current brakes. During braking, the metal wheels are exposed to a magnetic field from an electromagnet, generating eddy currents in the wheels. The eddy currents meet resistance as charges flow through the metal, thus dissipating energy as heat, and this acts to slow the wheels down. The faster the wheels are spinning, the stronger the effect, meaning that as the train slows the braking force is reduced, producing a smooth stopping motion.
Applications:-
In a fast varying magnetic field the induced currents, in good conductors, particularly copper and aluminium, exhibit diamagnetic-like repulsion effects on the magnetic field, and hence on the magnet and can create repulsive effects and even stable levitation, albeit with reasonably high power dissipation due to the high currents this entails.They can thus be used to induce a magnetic field in aluminum cans, which allows them to be separated easily from other recyclables. With a very strong handheld magnet, such as those made from neodymium, one can easily observe a very similar effect by rapidly sweeping the magnet over a coin with only a small separation. Depending on the strength of the magnet, identity of the coin, and separation between the magnet and coin, one may induce the coin to be pushed slightly ahead of the magnet - even if the coin contains no magnetic elements, such as the US penny.
Superconductors allow perfect, lossless conduction, which creates perpetually circulating eddy currents that are equal and opposite to the external magnetic field, thus allowing magnetic levitation. For the same reason, the magnetic field inside a superconducting medium will be exactly zero, regardless of the external applied field.In coin operated vending machines, eddy currents are used to detect counterfeit coins, or slugs. The coin rolls past a stationary magnet, and eddy currents slow its speed. The strength of the eddy currents, and thus the amount of slowing, depends on the conductivity of the coin's metal. Slugs are slowed to a different degree than genuine coins, and this is used to send them into the rejection slot.
Eddy currents are used in certain types of proximity sensors to observe the vibration and position of rotating shafts within their bearings. This technology was originally pioneered in the 1930s by researchers at General Electric using vacuum tube circuitry. In the late 1950s, solid-state versions were developed by Donald E. Bently at Bently Nevada Corporation. These sensors are extremely sensitive to very small displacements making them well suited to observe the minute vibrations (on the order of several thousandths of an inch) in modern turbomachinery. A typical proximity sensor used for vibration monitoring has a scale factor of 200 mV/mil. Widespread use of such sensors in turbomachinery has led to development of industry standards that prescribe their use and application. Examples of such standards are American Petroleum Institute (API) Standard 670 and ISO 7919.

Semiconductor

2009-11-01T09:52:00.000-08:00

A semiconductor is a material that has an electrical resistivity between that of a conductor and an insulator, that is, generally in the range 103 Siemens/cm to 10−8 S/cm. Devices made from semiconductor materials are the foundation of modern electronics, including radio, computers, telephones, and many other devices. Semiconductor devices include the various types of transistor, solar cells, many kinds of diodes including the light-emitting diode, the silicon controlled rectifier, and digital and analog integrated circuits. Solar photovoltaic panels are large semiconductor devices that directly convert light energy into electrical energy. An external electrical field may change a semiconductor's resistivity. In a metallic conductor, current is carried by the flow of electrons. In semiconductors, current can be carried either by the flow of electrons or by the flow of positively-charged "holes" in the electron structure of the material.
Common semiconducting materials are crystalline solids but amorphous and liquid semiconductors are known, such as mixtures of arsenic, selenium and tellurium in a variety of proportions. They share with better known semiconductors intermediate conductivity and a rapid variation of conductivity with temperature but lack the rigid crystalline structure of conventional semiconductors such as silicon and so are relatively insensitive to impurities and radiation damage.
Silicon is used to create most semiconductors commercially. Dozens of other materials are used, including germanium, gallium arsenide, and silicon carbide. A pure semiconductor is often called an “intrinsic” semiconductor. The conductivity, or ability to conduct, of common semiconductor materials can be drastically changed by adding other elements, called “impurities” to the melted intrinsic material and then allowing the melt to solidify into a new and different crystal. This process is called "doping".
Energy bands and electrical conduction:-
Like in other solids, the electrons in semiconductors can have energies only within certain bands (ie. ranges of levels of energy) between the energy of the ground state, corresponding to electrons tightly bound to the atomic nuclei of the material, and the free electron energy, which is the energy required for an electron to escape entirely from the material. The energy bands each correspond to a large number of discrete quantum states of the electrons, and most of the states with low energy (closer to the nucleus) are full, up to a particular band called the valence band. Semiconductors and insulators are distinguished from metals because the valence band in the semiconductor materials is very nearly full under usual operating conditions, thus causing more electrons to be available in the "conduction band," which is the band immediately above the valence band.
The ease with which electrons in a semiconductor can be excited from the valence band to the conduction band depends on the band gap between the bands, and it is the size of this energy bandgap that serves as an arbitrary dividing line (roughly 4 eV) between semiconductors and insulators.
In the picture of covalent bonds, an electron moves by hopping to a neighboring bond. Because of the Pauli exclusion principle it has to be lifted into the higher anti-bonding state of that bond. In the picture of delocalized states, for example in one dimension that is in a wire, for every energy there is a state with electrons flowing in one direction and one state for the electrons flowing in the other. For a net current to flow some more states for one direction than for the other direction have to be occupied and for this energy is needed. For a metal this can be a very small energy in the semiconductor the next higher states lie above the band gap. Often this is stated as: full bands do not contribute to the electrical conductivity. However, as the temperature of a semiconductor rises above absolute zero, there is more energy in the semiconductor to spend on lattice vibration and — more importantly for us — on lifting some electrons into an energy states of the conduction band. The current-carrying electrons in the conduction band are known as "free electrons", although they are often simply called "electrons" if context allows this usage to be clear.Electrons excited to the conduction band also leave behind electron holes, or unoccupied states in the valence band. Both the conduction band electrons and the valence band holes contribute to electrical conductivity. The holes themselves don't actually move, but a neighboring electron can move to fill the hole, leaving a hole at the place it has just come from, and in this way the holes appear to move, and the holes behave as if they were actual positively charged particles.
One covalent bond between neighboring atoms in the solid is ten times stronger than the binding of the single electron to the atom, so freeing the electron does not imply destruction of the crystal structure.
Holes: electron absence as a charge carrier:-
The motion of holes, which was introduced for semiconductors, can also be applied to metals, where the Fermi level lies within the conduction band. With most metals the Hall effect reveals electrons to be the charge carriers, but some metals have a mostly filled conduction band, and the Hall effect reveals positive charge carriers, which are not the ion-cores, but holes. Contrast this to some conductors like solutions of salts, or plasma. In the case of a metal, only a small amount of energy is needed for the electrons to find other unoccupied states to move into, and hence for current to flow. Sometimes even in this case it may be said that a hole was left behind, to explain why the electron does not fall back to lower energies: It cannot find a hole. In the end in both materials electron-phonon scattering and defects are the dominant causes for resistance.The energy distribution of the electrons determines which of the states are filled and which are empty. This distribution is described by Fermi-Dirac statistics. The distribution is characterized by the temperature of the electrons, and the Fermi energy or Fermi level. Under absolute zero conditions the Fermi energy can be thought of as the energy up to which available electron states are occupied. At higher temperatures, the Fermi energy is the energy at which the probability of a state being occupied has fallen to 0.5.The dependence of the electron energy distribution on temperature also explains why the conductivity of a semiconductor has a strong temperature dependency, as a semiconductor operating at lower temperatures will have fewer available free electrons and holes able to do the work.

Electric power

2009-10-24T10:59:00.000-07:00

Electric power is defined as the rate at which electrical energy is transferred by an electric circuit. The SI unit of power is the watt.When electric current flows in a circuit, it can transfer energy to do mechanical or thermodynamic work. Devices convert electrical energy into many useful forms, such as heat (electric heaters), light (light bulbs), motion (electric motors), sound (loudspeaker) or chemical changes. Electricity can be produced mechanically by generation, or chemically, or by direct conversion from light in photovoltaic cells, also it can be stored chemically in batteries.In alternating current circuits, energy storage elements such as inductance and capacitance may result in periodic reversals of the direction of energy flow. The portion of power flow that, averaged over a complete cycle of the AC waveform, results in net transfer of energy in one direction is known as real power (also referred to as active power). That portion of power flow due to stored energy, that returns to the source in each cycle, is known as reactive power.
The relationship between real power, reactive power and apparent power can be expressed by representing the quantities as vectors. Real power is represented as a horizontal vector and reactive power is represented as a vertical vector. The apparent power vector is the hypotenuse of a right triangle formed by connecting the real and reactive power vectors. This representation is often called the power triangle. Using the Pythagorean Theorem, the relationship among real, reactive and apparent power is:
apparent power)2 = (real power)2 + (reactive power)2
Real and reactive powers can also be calculated directly from the apparent power, when the current and voltage are both sinusoids with a known phase angle between them:
(real power) = (apparent power) * cos(theta)
(reactive power) = (apparent power) * sin(theta)
The ratio of real power to apparent power is called power factor and is a number always between 0 and 1.
The above theory of reactive power and the power triangle is true only when both the voltage and current is strictly sinusoidal. Therefore is more or less abandoned for low voltage distribution applications where the current normally is rather distorted. It can still be used for high voltage tranmission applications and, with some care, for medium voltage applications where the current normally is less distorted.

Earthing

2009-10-13T07:13:00.000-07:00

In electrical engineering, ground or earth may be the reference point in an electrical circuit from which other voltages are measured, or a common return path for electric current, or a direct physical connection to the Earth.Electrical circuits may be connected to ground (earth) for several reasons. In mains powered equipment, exposed metal parts are connected to ground to prevent contact with a dangerous voltage if electrical insulation fails. A connection to ground limits the voltage built up between power circuits and the earth, protecting circuit insulation from damage due to excessive voltage. Connections to ground limits the build-up of static electricity when handling flammable products or when repairing electronic devices. In some telegraph and power transmission circuits, the earth itself can be used as one conductor of the circuit, saving the cost of installing a separate return conductor.
For measurement purposes, the Earth serves as a (reasonably) constant potential reference against which other potentials can be measured. An electrical ground system should have an appropriate current-carrying capability in order to serve as an adequate zero-voltage reference level. In electronic circuit theory, a "ground" is usually idealized as an infinite source or sink for charge, which can absorb an unlimited amount of current without changing its potential. Where a ground connection has a significant resistance, the approximation of zero potential is no longer valid. Stray voltages or earth potential rise effects will occur, which may create noise in signals or if large enough will produce an electric shock hazard.
The use of the term ground (or earth) is so common in electrical and electronics applications that circuits in vehicles such as ships, aircraft, and spacecraft may be spoken of as having a "ground" connection without any actual connection to the Earth.
History:-

Long-distance electromagnetic telegraph systems from 1820 onwards used two or more wires to carry the signal and return currents. It was then discovered, probably by the German scientist Carl August Steinheil in 1836-1837 , that the ground could be used as the return path to complete the circuit, making the return wire unnecessary. However, there were problems with this system, exemplified by the transcontinental telegraph line constructed in 1861 by the Western Union Company between Saint Joseph, Missouri, and Sacramento, California. During dry weather, the ground connection often developed a high resistance, requiring water to be poured on the ground rod to enable the telegraph to work or phones to ring.
Later, when telephony began to replace telegraphy, it was found that the currents in the earth induced by power systems, electrical railways, other telephone and telegraph circuits, and natural sources including lightning caused unacceptable interference to the audio signals, and the two-wire system was reintroduced.

Refrigerator

2009-10-10T05:27:00.000-07:00

A refrigerator (often called a "fridge" for short) is a cooling appliance comprising a thermally insulated compartment and a heat pump—chemical or mechanical means—to transfer heat from it to the external environment, cooling the contents to a temperature below ambient. Refrigerators are extensively used to store foods which spoil from bacterial growth if not refrigerated. A device described as a "refrigerator" maintains a temperature a few degrees above the freezing point of water; a similar device which maintains a temperature below the freezing point of water is called a "freezer." The refrigerator is a relatively modern invention among kitchen appliances. It replaced the icebox, which had been a common household appliance for almost a century and a half prior. For this reason, a refrigerator is sometimes referred to as an "icebox".
History:-

Before the invention of the refrigerator, icehouses were used to provide cool storage for most of the year. Placed near freshwater lakes or packed with snow and ice during the winter, they were once very common. Using natural means to cool foods is still done today. On mountainsides, runoff from melting snow higher up is a convenient way to cool drinks, and during the winter months simply placing milk outside is sufficient to greatly extend its useful life.
In the 11th century, the Persian physicist and chemist, Ibn Sina (Avicenna), invented the refrigerated coil, which condenses aromatic vapours.This was a breakthrough in distillation technology and he made use of it in his steam distillation process, which requires refrigerated tubing, to produce essential oils.
The first known artificial refrigeration was demonstrated by William Cullen at the University of Glasgow in 1748. Between 1805, when Oliver Evans designed the first refrigeration machine that used vapor instead of liquid, and 1902 when Willis Haviland Carrier demonstrated the first air conditioner, scores of inventors contributed many small advances in cooling machinery. In 1850 or 1851, Dr. John Gorrie demonstrated an ice maker. In 1857, Australian James Harrison introduced vapor-compression refrigeration to the brewing and meat packing industries. Ferdinand Carré of France developed a somewhat more complex system in 1859. Unlike earlier compression-compression machines, which used air as a coolant, Carré's equipment contained rapidly expanding ammonia. The absorption refrigerator was invented by Baltzar von Platen and Carl Munters in 1922, while they were still students at the Royal Institute of Technology in Stockholm, Sweden. It became a worldwide success and was commercialized by Electrolux. Other pioneers included Charles Tellier, David Boyle, and Raoul Pictet. Carl von Linde was the first to patent and make a practical , and compact refrigerator.In a few exceptional cases, mechanical refrigeration systems had been adapted by the start of the 20th century for use in the homes of the very wealthy, and might be used for cooling both living and food storage areas. One early system was installed at the mansion of Walter Pierce, an oil company executive.General Electric sought to develop refrigerators of its own, and in 1915 the first Guardian unit was assembled in a back yard wash house as a predecessor to the Frigidaire. In 1916 Kelvinator and Servel introduced two units among a field of competing models. This number increased to 200 by 1920. In 1918, Kelvinator had a model with automatic controls.These home units usually required the installation of the mechanical parts, motor and compressor, in the basement or an adjacent room while the cold box was located in the kitchen. There was a 1922 model that consisted of a wooden cold box, water-cooled compressor, an ice cube tray and a 9 cubic foot compartment for $714. (A 1922 Model-T Ford cost about $450.) In 1923 Frigidaire introduced the first self-contained unit. About this same time porcelain covered metal cabinets began to appear. Ice cube trays were introduced more and more during the 1920s; up to this time freezing was not an auxiliary function of the modern refrigerator.The introduction of Freon expanded the refrigerator market during the 1930s. Separate freezers became common during the 1940s, the popular term at the time for the unit was a "deep freeze". But home units of these devices or "appliances" did not go into mass production until after WWII. The 1950s and 1960s saw technical advances like automatic defrosting and automatic ice making. Developments of the 1970s and 80s brought about more efficient refrigerators, even though environmental issues led to the banning of very effective (Freon) refrigerants. Early refrigerator models (1916 and on) featured a cold compartment for ice cube trays. Successful processing of fresh vegetables through freezing began in the late 1920s by the Postum Company (the forerunner of General Foods) which had acquired the technology when it bought the rights to Clarence Birdseye's successful fresh freezing methods.The first successful example of the benefits of frozen foods occurred when General Foods heiress Marjorie Merriweather Post (then wife of Joseph E. Davies, United States Ambassador to the Soviet Union) deployed commercial-grade freezers to Spaso House, the US Embassy in Moscow in advance of the Davies’ arrival. Post, fearful of the food processing safety observed in the USSR, then fully stocked the freezers with products from General Foods' Birdseye unit. The frozen food stores allowed the Davies to lavishly entertain and serve fresh frozen foods that would otherwise be out of season. Upon returning from Moscow, Post (who resumed her maiden name after divorcing Davies) directed General Foods to market frozen product to upscale restaurants.Introduction of home freezers as separate compartments (larger than necessary just for ice cubes), or as separate units, occurred in the United States in 1940. Frozen foods then began to make the transition from luxury to commonplace.

Electrical fan

2009-10-10T05:23:00.000-07:00

A electrical fan is an electrically powered device used to produce an airflow for the purpose of creature comfort (particularly in the heat), ventilation, exhaust, cooling or any other gaseous transport.Mechanically, a fan can be any revolving vane or vanes used for producing currents of air. Fans produce air flows with high volume and low pressure, as opposed to a gas compressor which produces high pressures at a comparatively low volume. A fan blade will often rotate when exposed to an air stream, and devices that take advantage of this, such as anemometers and wind turbines often have designs similar to that of a fan.In addition to their utilitarian function, vintage or antique fans, and in particular electric fans manufactured from the late 19th century through the 1950s, have become a recognized collectible category, and in the U.S.A. an active collector club, the Antique Fan Collectors Association, supports the hobby.New to the market are sleek portable fans that showcase a modern design sensibility. The New York Times lamented that inexpensive and effective fans abound at drug and discount stores, but they are often eyesores. The writer quoted contemporary ceiling fan designer Ron Rezek as saying: “Portable fans are the ugly ducklings of the fan industry. Not many designers, including myself, have tackled them[2].” Rezek praised several appealing contemporary fan designs that have found alternatives to the traditional metal cage and have incorporated innovative approaches to safety, such as the Otto fan by Swiss designer Carlo Borer.
History:-

The Industrial Revolution in the late 19th century introduced belt-driven fans powered by factory waterwheels. Attaching wooden or metal blades to shafts overhead that were used to drive the machinery, the first industrial fans were developed. One of the first workable mechanical fans was built by Omar-Rajeen Jumala in 1832. He called his invention, a kind of a centrifugal fan, an Air Pump. Centrifugal fans were successfully tested inside coal mines and factories in 1832-1834. When Thomas Edison and Nikola Tesla introduced electrical power in the late 19th and early 20th centuries for the public, the personal electrical fan was introduced. Between 1882 and 1886, Dr. Schuyler Skaats Wheeler developed the two-bladed desk fan, a type of personal electric fan. It was commercially marketed by the American firm Crocker & Curtis electric motor company. In 1882, Philip Diehl introduced the electric ceiling fan. Diehl is considered the father of the modern electric fan. In the late 19th century, electric fans were used only in commercial establishments or in well-to-do households. Heat-convection fans fueled by alcohol, oil, or kerosene were common around the turn of the 20th century.The first American fans were made from around the late 1890s to the early 1920s. They had brass blades, a lot of them also had brass cages, and though they were built very well internally, they were far from finger safe, as a lot of them had cage openings so big that one could put an entire hand or arm right through it. Many children had hands and fingers severely injured by those fans.

Electrical iron

2009-10-10T05:17:00.000-07:00

An electrical iron is a small appliance used in ironing to remove wrinkles from fabric. Ironing works by loosening the ties between the long chains of molecules that exist in polymer fiber materials. With the heat and the weight of the ironing plate, the fibers are stretched and the fabric maintains its new shape when cool. Some materials such as cotton require the use of water to loosen the intermolecular bonds. Many materials developed in the twentieth century are advertised as needing little or no ironing.
History:-

Metal pans filled with charcoal were used for smoothing fabrics in China in the 1st century BC. From the 17th century, sadirons or sad irons began to be used. They were thick slabs of cast iron, delta-shaped and with a handle, heated in a fire. These were also called flat irons. A later design consisted of an iron box which could be filled with hot coals, which had to be periodically aerated by attaching a bellows. In Kerala in India, burning coconut shells were used instead of charcoal, as they have a similar heating capacity. This method is still in use as a backup device since power outages are frequent. Other box irons had heated metal inserts instead of hot coals. Another solution was a cluster of solid irons that were heated from the single source: as the iron currently in use cools down, it can be quickly replaced by another one that is hot.In the industrialized world, these designs have been superseded by the electric iron, which uses resistive heating from an electric current. The hot plate, called the sole plate, is made of aluminium or stainless steel. The heating element is controlled by a thermostat which switches the current on and off to maintain the selected temperature. The invention of the resistively heated electric iron is credited to Henry W. Seely of New York in 1882. In the same year an iron heated by a carbon arc was introduced in France, but was too dangerous to be successful. The early electric irons had no easy way to control their temperature, and the first thermostatically controlled electric iron appeared in the 1920s. Later, steam was used to iron clothing. Credit for the invention of the steam iron goes to Thomas Sears.

Multimeter

2009-10-10T05:13:00.000-07:00

A multimeter or a multitester, also known as a volt/ohm meter or VOM, is an electronic measuring instrument that combines several measurement functions in one unit. A typical multimeter may include features such as the ability to measure voltage, current and resistance. There are two categories of multimeters, analog multimeters and digital multimeters.A multimeter can be a hand-held device useful for basic fault finding and field service work or a bench instrument which can measure to a very high degree of accuracy. They can be used to troubleshoot electrical problems in a wide array of industrial and household devices such as batteries, motor controls, appliances, power supplies, and wiring systems.Multimeters are available in a wide ranges of features and prices. Cheap multimeters can cost less than US$10, while the top of the line multimeters can cost more than US$5000.
History:-

Scientists originally used galvanometers to measure current. A galvanometer may be wired to measure resistance (given a known voltage source) or voltage (given a fixed resistance). While appropriate for primitive lab use, switching from one setup to another is inconvenient in the field.Multimeters were invented in the early 1920s as radio receivers and other vacuum tube electronic devices became more common. The invention of the first multimeter is attributed to Post Office engineer Donald Macadie, who became dissatisfied with having to carry many separate instruments required for the maintenance of the telecommunication circuits.Macadie took his idea to the Automatic Coil Winder and Electrical Equipment Company.The first AVO was put on sale in 1923, and although it was initially a DC-only instrument many of its features remained almost unaltered right through to the last Model 8.As modern systems become more complicated, the multimeter is becoming more complex or may be supplemented by more specialized equipment in a technician's toolkit. For example, where a general-purpose multimeter might only test for short-circuits, conductor resistance and some coarse measure of insulation quality, a modern technician may use a hand-held analyzer to test several parameters in order to validate the performance of a network cable.

Rectifier

2009-10-10T05:00:00.000-07:00

A rectifier is an electrical device that converts alternating current (AC) to direct current (DC), a process known as rectification. Rectifiers have many uses including as components of power supplies and as detectors of radio signals. Rectifiers may be made of solid state diodes, vacuum tube diodes, mercury arc valves, and other components.A device which performs the opposite function (converting DC to AC) is known as an inverter.When only one diode is used to rectify AC (by blocking the negative or positive portion of the waveform), the difference between the term diode and the term rectifier is merely one of usage, i.e., the term rectifier describes a diode that is being used to convert AC to DC. Almost all rectifiers comprise a number of diodes in a specific arrangement for more efficiently converting AC to DC than is possible with only one diode. Before the development of silicon semiconductor rectifiers, vacuum tube diodes and copper(I) oxide or selenium rectifier stacks were used.
Half wave rectification:-

In half wave rectification, either the positive or negative half of the AC wave is passed, while the other half is blocked. Because only one half of the input waveform reaches the output, it is very inefficient if used for power transfer. Half-wave rectification can be achieved with a single diode in a one-phase supply, or with three diodes in a three-phase supply.
Full wave rectification:-

A full-wave rectifier converts the whole of the input waveform to one of constant polarity (positive or negative) at its output. Full-wave rectification converts both polarities of the input waveform to DC (direct current), and is more efficient. However, in a circuit with a non-center tapped transformer, four diodes are required instead of the one needed for half-wave rectification. (See semiconductors, diode). Four rectifiers arranged this way are called a diode bridge or bridge rectifier.
Applicatinos:-

The primary application of rectifiers is to derive DC power from an AC supply. Virtually all electronic devices require DC, so rectifiers find uses inside the power supplies of virtually all electronic equipment.Converting DC power from one voltage to another is much more complicated. One method of DC-to-DC conversion first converts power to AC (using a device called an inverter), then use a transformer to change the voltage, and finally rectifies power back to DC.Rectifiers also find a use in detection of amplitude modulated radio signals. The signal may or may not be amplified before detection but if un-amplified a very low voltage drop diode must be used. When using a rectifier for demodulation the capacitor and load resistance must be carefully matched. Too low a capacitance will result in the high frequency carrier passing to the output and too high will result in the capacitor just charging and staying charged.

Electrical generator

2009-10-10T04:56:00.000-07:00

In electricity generation, an electrical generator is a device that converts mechanical energy to electrical energy. The reverse conversion of electrical energy into mechanical energy is done by a motor; motors and generators have many similarities. A generator forces electric charges to move through an external electrical circuit, but it does not create electricity or charge, which is already present in the wire of its windings. It is somewhat analogous to a water pump, which creates a flow of water but does not create the water inside. The source of mechanical energy may be a reciprocating or turbine steam engine, water falling through a turbine or waterwheel, an internal combustion engine, a wind turbine, a hand crank, compressed air or any other source of mechanical energy.
History:-

Before the connection between magnetism and electricity was discovered, electrostatic generators were invented that used electrostatic principles. These generated very high voltages and low currents. They operated by using moving electrically charged belts, plates and disks to carry charge to a high potential electrode.

Induction motor

2009-10-10T04:42:00.000-07:00

An induction motor is a type of alternating current motor where power is supplied to the rotating device by means of electromagnetic induction. It is also called asynchronous motor.An electric motor converts electrical power to mechanical power in its rotor (rotating part). There are several ways to supply power to the rotor. In a DC motor this power is supplied to the armature directly from a DC source, while in an induction motor this power is induced in the rotating device. An induction motor is sometimes called a rotating transformer because the stator (stationary part) is essentially the primary side of the transformer and the rotor (rotating part) is the secondary side. The primary side's currents evokes a magnetic field which interacts with the secondary sides mmf to produce a resultant torque, henceforth serving the purpose of producing mechanical energy. Induction motors are widely used, especially polyphase induction motors, which are frequently used in industrial drives.
History:-

The induction motor with a wrapped rotor was invented by Nikola Tesla in 1882 in France but the initial patent was issued in 1888 after Tesla had moved to the United States. In his scientific work, Tesla laid the foundations for understanding the way the motor operates. The induction motor with a cage was invented by Mikhail Dolivo-Dobrovolsky about a year later in Europe. Technological development in the field has improved to where a 100 hp (74.6 kW) motor from 1976 takes the same volume as a 7.5 hp (5.5 kW) motor did in 1897. Currently, the most common induction motor is the cage rotor motor. Principle of operation and comparison to synchronous motors:-The basic difference between an induction motor and a synchronous AC motor is that in the latter a current is supplied onto the rotor. This then creates a magnetic field which, through magnetic interaction, links to the rotating magnetic field in the stator which in turn causes the rotor to turn. It is called synchronous because at steady state the speed of the rotor is the same as the speed of the rotating magnetic field in the stator.
By way of contrast, the induction motor does not have any direct supply onto the rotor; instead, a secondary current is induced in the rotor. To achieve this, stator windings are arranged around the rotor so that when energised with a polyphase supply they create a rotating magnetic field pattern which sweeps past the rotor. This changing magnetic field pattern induces current in the rotor conductors. These currents interact with the rotating magnetic field created by the stator and in effect causes a rotational motion on the rotor.
However, for these currents to be induced, the speed of the physical rotor must be less than the speed of the rotating magnetic field in the stator, or else the magnetic field will not be moving relative to the rotor conductors and no currents will be induced. If by some chance this happens, the rotor typically slows slightly until a current is re-induced and then the rotor continues as before. This difference between the speed of the rotor and speed of the rotating magnetic field in the stator is called slip. It is unitless and is the ratio between the relative speed of the magnetic field as seen by the rotor (the slip speed) to the speed of the rotating stator field. Due to this an induction motor is sometimes referred to as an asynchronous machine.

Construction:-
The stator consists of wound 'poles' that carry the supply current to induce a magnetic field that penetrates the rotor. In a very simple motor, there would be a single projecting piece of the stator (a salient pole) for each pole, with windings around it; in fact, to optimize the distribution of the magnetic field, the windings are distributed in many slots located around the stator, but the magnetic field still has the same number of north-south alternations.Induction motors are most commonly built to run on single-phase or three-phase power, but two-phase motors also exist. In theory, two-phase and more than three phase induction motors are possible; many single-phase motors having two windings and requiring a capacitor can actually be viewed as two-phase motors, since the capacitor generates a second power phase 90 degrees from the single-phase supply and feeds it to a separate motor winding. Single-phase power is more widely available in residential buildings, but cannot produce a rotating field in the motor (the field merely oscillates back and forth), so single-phase induction motors must incorporate some kind of starting mechanism to produce a rotating field. They would, using the simplified analogy of salient poles, have one salient pole per pole number; a four-pole motor would have four salient poles. Three-phase motors have three salient poles per pole number, so a four-pole motor would have twelve salient poles. This allows the motor to produce a rotating field, allowing the motor to start with no extra equipment and run more efficiently than a similar single-phase motor.

Single phase transformer

2009-10-10T04:25:00.000-07:00

In electrical engineering, single-phase electric power refers to the distribution of alternating current electric power using a system in which all the voltages of the supply vary in unison. Single-phase distribution is used when loads are mostly lighting and heating, with few large electric motors. A single-phase supply connected to an alternating current electric motor does not produce a revolving magnetic field; single-phase motors need additional circuits for starting, and such motors are uncommon above 10 or 20 kW in rating.In contrast, in a three-phase system, the currents in each conductor reach their peak instantaneous values sequentially, not simultaneously; in each cycle of the power frequency, first one, then the second, then the third current reaches its maximum value. The waveforms of the three supply conductors are offset from one another in time by one-third of their period.
Splitting out:-

No arrangement of transformers can convert a single-phase load into a balanced load on a polyphase system. A single-phase load may be powered from a three-phase distribution system either by connection between a phase and neutral or by connecting the load between two phases. The load device must be designed for the voltage in each case. The neutral point in a three phase system exists at the mathematical center of an equilateral triangle formed by the three phase points, and the phase-to-phase voltage is accordingly times the phase-to-neutral voltage.
In North America, a typical three-phase system will have 208 volts between the phases and 120 volts between phase and neutral. If heating equipment designed for the 240-volt three-wire single phase system is connected to two phases of a 208 volt supply, it will only produce 75% of its rated heating effect. Single-phase motors may have taps to allow their use on either 208 V or 240 V supplies.
On higher voltage systems (KV) where a single phase transformer is in use to supply a low voltage system the method of splitting varies. In North America utility distribution practice, the primary of the step-down transformer is wired across a single high voltage feed wire and neutral, at least for smaller supplies. Rural distribution may be a single phase at a medium voltage; in some areas single wire earth return distribution is used when customers are very far apart. In Britain the step-down primary is wired phase-phase.
Applications:-

Single-phase power distribution is widely used especially in rural areas, where the cost of a three-phase distribution network is high and motor loads are small and uncommon.
High power systems, say, hundreds of kVA or larger, are nearly always three phase. The largest supply normally available as single phase varies according to the standards of the electrical utility. In the UK a single-phase household supply may be rated 100 A or even 125 A, meaning that there is little need for 3 phase in a domestic or small commercial environment. Much of the rest of Europe has traditionally had much smaller limits on the size of single phase supplies resulting in even houses being supplied with 3 phase in urban areas with three-phase supply networks.
In North America, individual residences and small commercial buildings with services up to about 100 kV·A will usually have three-wire single-phase distribution, often with only one customer per distribution transformer. In exceptional cases larger single-phase three-wire services can be provided, usually only in remote areas where polyphase distribution is not available. In rural areas farmers who wish to use three-phase motors may install a phase converter if only a single-phase supply is available. Larger consumers such as large buildings, shopping centres, factories, office blocks, and multiple-unit apartment blocks will have three-phase service. In densely-populated areas of cities, network power distribution is used with many customers and many supply transformers connected to provide hundreds or thousands of kV·A load concentrated over a few hundred square metres.

Grownd wiring

2009-10-10T03:27:00.000-07:00

Electrical wiring

2009-10-10T03:25:00.000-07:00

Electrical wiring in general refers to insulated conductors used to carry electricity, and associated devices. now i shall explan that, electrical wiring as used to provide power in buildings and structures, commonly referred to as building wiring. This article is intended to describe common features of electrical wiring that should apply worldwide.the peroson who do the work this is called electrician.

Cable

2009-09-11T06:04:00.000-07:00

A cable is two or more wires or ropes running side by side and bonded, twisted or braided together to form a single assembly. In mechanics, cables are used for lifting and hauling; in electricity they are used to carry electrical currents. An optical cable contains one or more optical fibers in a protective jacket that supports the fibers. Mechanical cable is more specifically called wire rope.
History:-
Ropes made of multiple strands of natural fibers such as hemp, sisal, manila, and cotton have been used for millennia for hoisting and hauling. By the 19th century, deepening of mines and construction of large ships increased demand for stronger cables. Invention of improved steelmaking techniques made high quality steel available at lower cost, and so wire ropes became common in mining and other industrial applications. By the middle of the 19th century, manufacture of large submarine telegraph cables was done using machiners similar to that used for manufacture of mechanical cables.
In the 19th century and early 20th century, electrical cable was often insulated using cloth, rubber and paper. Plastic materials are generally used today, except for high reliability power cables.
Electrical cable:-
Electrical cables may be made more flexible by stranding the wires. In this process, smaller individual wires are twisted or braided together to produce larger wires that are more flexible than solid wires of similar size. Bunching small wires before concentric stranding adds the most flexibility. Copper wires in a cable may be bare, or they may be coated with a thin layer of another material: most often tin but sometimes gold, silver or some other material. Tin, gold, and silver are much less prone to oxidisation than copper, which may lengthen wire life, and makes soldering easier. Tight lays during stranding makes the cable extensible (CBA - as in telephone handset cords).
Cables can be securely fastened and organized, such as by using cable trees with the aid of cable ties or cable lacing. Continuous-flex or flexible cables used in moving applications within cable carriers can be secured using strain relief devices or cable ties. Copper corrodes easily and so should be layered with Lacquer.
At high frequencies, current tends to run along the surface of the conductor and avoid the core. This is known as the skin effect. It may change the relative desirability of solid versus stranded wires.

Cartridge type fuse

2009-09-08T21:43:00.000-07:00

In electronics and electrical engineering a fuse (short for fusible link) is a type of sacrificial overcurrent protection device. Its essential component is a metal wire or strip that melts when too much current flows, which interrupts the circuit in which it is connected. Short circuit, overload or device failure is often the reason for excessive current.
A fuse interrupts excessive current (blows) so that further damage by overheating or fire is prevented. Wiring regulations often define a maximum fuse current rating for particular circuits. Overcurrent protection devices are essential in electrical systems to limit threats to human life and property damage. Fuses are selected to allow passage of normal current and of excessive current only for short periods.
A fuse was patented by Thomas Edison in 1890 [1] as part of his successful electric distribution system.
Opration:-
A fuse consists of a metal strip or wire fuse element, of small cross-section compared to the circuit conductors, mounted between a pair of electrical terminals, and (usually) enclosed by a non-conducting and non-combustible housing. The fuse is arranged in series to carry all the current passing through the protected circuit. The resistance of the element generates heat due to the current flow. The size and construction of the element is (empirically) determined so that the heat produced for a normal current does not cause the element to attain a high temperature. If too high a current flows, the element rises to a higher temperature and either directly melts, or else melts a soldered joint within the fuse, opening the circuit.
When the metal conductor parts, an electric arc forms between the un-melted ends of the element. The arc grows in length until the voltage required to sustain the arc is higher than the available voltage in the circuit, terminating current flow. In alternating current circuits the current naturally reverses direction on each cycle, greatly enhancing the speed of fuse interruption. In the case of a current-limiting fuse, the arc voltage builds up quickly enough to essentially stop the fault current before the first peak of the ac waveform. This effect significantly limits damage to downstream protected devices.
The fuse element is made of zinc, copper, silver, aluminum, or alloys to provide stable and predictable characteristics. The fuse ideally would carry its rated current indefinitely, and melt quickly on a small excess. The element must not be damaged by minor harmless surges of current, and must not oxidize or change its behavior after possibly years of service.
The fuse elements may be shaped to increase heating effect. In large fuses, current may be divided between multiple strips of metal. A dual-element fuse may contain a metal strip that melts instantly on a short-circuit, and also contain a low-melting solder joint that responds to long-term overload of low values compared to a short-circuit. Fuse elements may be supported by steel or nichrome wires, so that no strain is placed on the element, but a spring may be included to increase the speed of parting of the element fragments.
The fuse element may be surrounded by air, or by materials intended to speed the quenching of the arc. Silica sand or non-conducting liquids may be used.
Speed:-
The speed at which a fuse blows depends on how much current flows through it and the material of which the fuse is made. The operating time is not a fixed interval, but decreases as the current increases. Fuses have different characteristics of operating time compared to current, characterized as "fast-blow", "slow-blow" or "time-delay", according to time required to respond to an overcurrent condition. A standard fuse may require twice its rated current to open in one second, a fast-blow fuse may require twice its rated current to blow in 0.1 seconds, and a slow-blow fuse may require twice its rated current for tens of seconds to blow.
Fuse selection depends on the load's characteristics. Semiconductor devices may use a fast or ultrafast fuse since semiconductor devices heat rapidly when excess current flows. The fastest blowing fuses are designed for the most sensitive electrical equipment, where even a short exposure to an overload current could be very damaging. Normal fast-blow fuses are the most general purpose fuses. The time delay fuse (also known as anti-surge, or slow-blow) are designed to allow a current which is above the rated value of the fuse to flow for a short period of time without the fuse blowing. These types of fuse are used on equipment such as motors, which draw a large initial current for a few milliseconds after they have been switched on.
Rated voltage:-
Voltage rating of the fuse must be greater than or equal to what would become the open circuit voltage. For example, a glass tube fuse rated at 32 volts would not reliably interrupt current from a voltage source of 120 or 230 V. If a 32 V fuse attempts to interrupt the 120 or 230 V source, an arc may result. Plasma inside that glass tube fuse may continue to conduct current until current eventually so diminishes that plasma reverts to an insulating gas. Rated voltage should be larger than the maximum voltage source it would have to disconnect. This requirement applies to every type of fuse.
Rated voltage remains same for any one fuse, even when similar fuses are connected in series. Connecting fuses in series does not increase the rated voltage.
Medium-voltage fuses rated for a few thousand volts are never used on low voltage circuits, because of their cost and because they cannot properly clear the circuit when operating at very low voltages.
Voltage drop:-
A voltage drop across the fuse is usually provided by its manufacturer. Resistance may change when a fuse becomes hot due to energy dissipation while conducting higher currents. This resulting voltage drop should be taken into account, particularly when using a fuse in low-voltage applications. Voltage drop often is not significant in more traditional wire type fuses, but can be significant in other technologies such as resettable fuse (PPTC) type fuses.

Fuse material

2009-09-08T11:23:00.000-07:00

Photo diode

2009-08-25T00:55:00.000-07:00

A photodiode is a type of photodetector capable of converting light into either current or voltage, depending upon the mode of operation.
Photodiodes are similar to regular semiconductor diodes except that they may be either exposed (to detect vacuum UV or X-rays) or packaged with a window or optical fiber connection to allow light to reach the sensitive part of the device. Many diodes designed for use specifically as a photodiode will also use a PIN junction rather than the typical PN junction.
Polarity:-
Some photodiodes will look like the picture to the right, that is, similar to a light emitting diode. They will have two leads, or wires, coming from the bottom. The shorter end of the two is the cathode, while the longer end is the anode. See below for a schematic drawing of the anode and cathode side. Under forward bias, conventional current will pass from the anode to the cathode, following the arrow in the symbol. Photocurrent flows in the opposite direction.
Principal of operation:-
A photodiode is a PN junction or PIN structure. When a photon of sufficient energy strikes the diode, it excites an electron, thereby creating a mobile electron and a positively charged electron hole. If the absorption occurs in the junction's depletion region, or one diffusion length away from it, these carriers are swept from the junction by the built-in field of the depletion region. Thus holes move toward the anode, and electrons toward the cathode, and a photocurrent is produced.

Zener diode

2009-08-25T00:46:00.000-07:00

A Zener diode is a type of diode that permits current in the forward direction like a normal diode, but also in the reverse direction if the voltage is larger than the breakdown voltage known as "Zener knee voltage" or "Zener voltage". The device was named after Clarence Zener, who discovered this electrical property. A conventional solid-state diode will not allow significant current if it is reverse-biased below its reverse breakdown voltage. When the reverse bias breakdown voltage is exceeded, a conventional diode is subject to high current due to avalanche breakdown. Unless this current is limited by external circuitry, the diode will be permanently damaged. In case of large forward bias (current in the direction of the arrow), the diode exhibits a voltage drop due to its junction built-in voltage and internal resistance. The amount of the voltage drop depends on the semiconductor material and the doping concentrations.
A Zener diode exhibits almost the same properties, except the device is specially designed so as to have a greatly reduced breakdown voltage, the so-called Zener voltage. A Zener diode contains a heavily doped p-n junction allowing electrons to tunnel from the valence band of the p-type material to the conduction band of the n-type material. In the atomic scale, this tunneling corresponds to the transport of valence band electrons into the empty conduction band states; as a result of the reduced barrier between these bands and high electric fields that are induced due to the relatively high levels of dopings on both sides. A reverse-biased Zener diode will exhibit a controlled breakdown and allow the current to keep the voltage across the Zener diode at the Zener voltage. For example, a diode with a Zener breakdown voltage of 3.2 V will exhibit a voltage drop of 3.2 V if reverse bias voltage applied across it is more than its Zener voltage. However, the current is not unlimited, so the Zener diode is typically used to generate a reference voltage for an amplifier stage, or as a voltage stabilizer for low-current applications.
The breakdown voltage can be controlled quite accurately in the doping process. While tolerances within 0.05% are available, the most widely used tolerances are 5% and 10%.
Another mechanism that produces a similar effect is the avalanche effect as in the avalanche diode. The two types of diode are in fact constructed the same way and both effects are present in diodes of this type. In silicon diodes up to about 5.6 volts, the Zener effect is the predominant effect and shows a marked negative temperature coefficient. Above 5.6 volts, the avalanche effect becomes predominant and exhibits a positive temperature coefficient.
Uses:-

Zener diodes are widely used to regulate the voltage across a circuit. When connected in parallel with a variable voltage source so that it is reverse biased, a Zener diode conducts when the voltage reaches the diode's reverse breakdown voltage. From that point it keeps the voltage at that value.

Pentode

2009-08-25T00:34:00.000-07:00

pentode is an electronic device having five active electrodes. The term most commonly applies to a three-grid vacuum tube, which was invented by the Dutchman Bernhard D.H. Tellegen in 1926.
Advantages over the tetrode:-

A tetrode could supply sufficient power to a speaker or transmitter, and offered a larger amplification factor than the earlier triode.
However, the positively charged screen grid can collect the secondary electrons emitted from the anode, which can cause increased current toward the screen grid, and cause the anode current to decrease with increasing anode voltage over part of the Ia/Va characteristic. Tellegen introduced an additional electrode, called the suppressor grid, which solved the problem of secondary emission. It does this by being held at a low potential, usually either grounded or connected to the cathode. The secondary emission still occurs, but the electrons can no longer reach the screen grid, since they have less energy than the primary electrons and hence cannot pass the grounded suppressor grid. Therefore these secondary electrons are re-collected by the anode.

Usage:-

Pentode valves were first used in consumer-type radio receivers. A well-known pentode type, the EF50, was designed before the start of the World War II, and was extensively used in radar sets and other military electronic equipment. The pentode contributed to the electronic preponderance of the Allies. After WW II, pentodes were widely used in TV receivers, particularly the successor to the EF50, the EF80. Vacuum tubes were replaced by transistors during the 1960s. However, they continue to be used in certain applications, including high-power radio transmitters and (because of their well-known valve sound) in high-end and professional audio applications, microphone preamplifiers and electric guitar amplifiers. Large stockpiles in countries of the former Soviet Union have provided a continuing supply of such devices, some designed for other purposes but adapted to audio use, such as the GU-50 transmitter tube.

Tetrode

2009-08-25T00:27:00.000-07:00

A tetrode is an electronic device having four active electrodes. The term most commonly applies to a two-grid vacuum tube. It has the three electrodes of a triode and an additional screen grid which significantly changes its behaviour.
Control grid:-
The grid nearest the cathode is the "control grid"; the voltage applied to it causes the anode current to vary. In normal operation, with a resistive load, this varying current will result in varying (AC) voltage measured at the anode. With proper biasing, this voltage will be an amplified (but inverted) version of the AC voltage applied to the control grid, thus the tetrode can provide voltage gain.
Screen grid:-
The second grid, called "screen grid" or sometimes "shield grid", provides a screening effect, isolating the control grid from the anode. This helps to suppress unwanted oscillation, and to reduce an undesirable effect in triodes called the "Miller effect", where the gain of the tube causes a feedback effect which increases the apparent capacitance of the tube's grid, limiting the tube's high-frequency gain. In normal operation the screen grid is connected to a positive voltage, and bypassed to the cathode with a capacitor. This shields the grid from the anode, reducing Miller capacitance between those two electrodes to a very low level and improving the tube's gain at high frequencies. When the tetrode was introduced, a typical triode had an input capacitance of about 5 pF, but the screen grid reduced this capacitance to about 0.01 pF.
As the screen grid is positively charged, it collects electrons, which causes current to flow in the screen grid circuit. This uses power and heats the screen grid; if the screen heats up enough it can melt and destroy the tube. There are two sources of electrons collected by the screen grid—in addition to the electrons emitted by the cathode, the screen grid can also collect secondary electrons ejected from the anode by the impact of the energetic primary electrons. Secondary emission can increase enough to decrease the anode current, since a single primary electron can eject more than one secondary electron. The reduction in anode current is because the external anode current (through the connection pin) is due to the cathode-to-anode current minus the secondary emission current. This can give the tetrode valve a distinctive negative resistance characteristic, sometimes called "tetrode kink". This is usually undesirable, although it can be exploited as in the dynatron oscillator. The secondary emission can be suppressed by adding a suppressor grid, making a pentode, or beam plates to make a beam tetrode/kinkless tetrode.
The positive influence of the screen grid in the vicinity of the control grid allows a designer to shift the control grid operating voltage range entirely into the negative region (a triode of similar geometry would likely require positive grid drive to attain the same maximum anode current). When any grid is driven positive relative to the cathode it can intercept electrons from the cathode, loading the drive circuitry. If the input signal causes the control grid to become positive (where current flow begins), nonlinearity is to be expected. (The control grid draws no current while negative—high impedance—but draws current while positive—low impedance.) With the control grid operating entirely in the negative region, and with the RF shielding afforded by the screen grid, tetrode input impedance is quite high even at high frequencies. Gain can be nearly flat from DC to full frequency. Linearity is good. Power gain in excess of 10,000 is possible.
The triode vacuum tube also develops a "space charge" between the cathode and control grid, which reduces its gain, especially at low anode voltages. The screen grid of a tetrode neutralizes the space charge and increases the tube's gain.
Power tetrodes are commonly used in radio transmitting equipment, because the need for neutralization is less than with triodes (see Radio transmitter design and Valve amplifier for more details). Screen current does represent loss. Some tube designers attempt to minimize screen current by placing each wire in the screen mesh directly behind a corresponding wire in the control grid mesh. Propagating electrons emerge from the control grid as a projected image of openings in the grid. By placing the screen in the shadow of the control grid, interception of electrons by the screen is minimized in normal operation. Screen current is negligible in many designs. Shadow grids are used in a variety of forms for a number of applications.
More than one screen grid can be used. For example the pentagrid converter has two. A tetrode can be converted to act as a triode by connecting the screen grid to the anode.
Circuit design considerations Under certain operating conditions, the tetrode exhibits negative resistance due to secondary emission of electrons from the anode (to the screen). The shape of the characteristic curve of a tetrode operated in this region led to the term "tetrode kink". In general, if the anode voltage exceeds the screen voltage, this region is avoided, and good performance can be expected. But this lower limit on total tube voltage drop prevents widespread adoption of tetrodes for consumer amplification applications. Secondary emissions from a screen have the effect of pulling the screen upward, toward the anode voltage. This implies the need for both source and sink current capability in the ideal screen power supply. A bleeder resistor can usually be selected to prevent the screen voltage from getting out of control. Arcs from the anode generally hit the screen. As such, special care is required in design of the socket wiring, to provide a direct discharge path for arc current. The undesirable nature of the tetrode kink led tube designers to add a third grid, called the suppressor grid; the resulting vacuum tube is called a pentode. More modern tubes have anodes treated to minimise secondary emission.
The negative resistance operating region of the tetrode is exploited in the dynatron oscillator, although this was practical only with earlier tubes with high secondary emission.
Invention:-
The tetrode tube was developed by Dr. Walter H. Schottky of Siemens & Halske GMBH in Germany during World War I. Thousands of variations of the tetrode design, as well as its later development the pentode, have been manufactured since then, although vacuum tubes in low-power equipment have been almost totally superseded by solid-state semiconductor devices.

Triode

2009-08-25T00:05:00.000-07:00

A triode is an electronic amplification device having three active electrodes. The term most commonly applies to a vacuum tube (or valve in British English) with three elements: the filament or cathode, the grid, and the plate or anode. The triode vacuum tube is the first electronic amplification device.
Invention:-

The original three-element device was patented in 1908 by Lee De Forest who developed it from his original two-element 1906 Audion. The Audion did provide amplification. However it was not until around 1912 that other researchers, while attempting to improve the service life of the audion, stumbled on the principle of the true vacuum tube. The name triode appeared later, when it became necessary to distinguish it from other generic kinds of vacuum tubes with more or less elements (eg diodes, tetrodes, pentodes etc.). The Audion tubes deliberately contained some gas at low pressure. The name triode is only applied to vacuum tubes which have been evacuated of as much gas as possible. There was a parallel invention of the triode in charge of Austrian Robert von Lieben.
Operation:-
The principle of its operation is that, as with a thermionic diode, the heated cathode(either directly or indirectly by means of a filament) causes a space charge of electrons that may be attracted to the positively charged plate (anode in UK parlance) and create a current. Applying a negative charge to the control grid will tend to repel some of the (also negatively charged) electrons back towards the cathode: the larger the charge on the grid, the smaller the current to the plate. If an AC signal is superimposed on the DC bias of the grid, an amplified version of the AC signal appears in the plate circuit.

Applications:-

Although triodes are now largely obsolete in consumer electronics, having been replaced by the transistor, triodes continue to be used in certain high-end and professional audio applications, as well as in microphone preamplifiers and electric guitar amplifiers. Some guitarists routinely drive their amplifiers to the point of saturation, in order to produce a desired distortion tone. Many people prefer the sound of triodes in such an application, since the distortion of a tube amplifier, which has a "soft" saturation characteristic, can be more pleasing to the ear than that of a typical solid-state amplifier, which is linear up to the limits of its supply voltage and then clips abruptly. However, this typically only applied to the power stage of a tube amplifier.

Diode

2009-08-24T23:53:00.000-07:00

In electronics, a diode is a two-terminal device (thermionic diodes may also have one or two ancillary terminals for a heater).
Diodes have two active electrodes between which the signal of interest may flow, and most are used for their unidirectional electric current property.
The unidirectionality most diodes exhibit is sometimes generically called the rectifying property. The most common function of a diode is to allow an electric current in one direction (called the forward biased condition) and to block the current in the opposite direction (the reverse biased condition). Thus, the diode can be thought of as an electronic version of a check valve.
Real diodes do not display such a perfect on-off directionality but have a more complex non-linear electrical characteristic, which depends on the particular type of diode technology. Diodes also have many other functions in which they are not designed to operate in this on-off manner.
Early diodes included “cat’s whisker” crystals and vacuum tube devices (also called thermionic valves). Today most diodes are made of silicon, but other semiconductors such as germanium are sometimes used.

History:-
Although the crystal (solid state) diode was popularized before the thermionic diode, thermionic and solid state diodes were developed in parallel.
The basic principle of operation of thermionic diodes was discovered by Frederick Guthrie in 1873.[1] Guthrie discovered that a positively-charged electroscope could be discharged by bringing a grounded piece of white-hot metal close to it (but not actually touching it). The same did not apply to a negatively charged electroscope, indicating that the current flow was only possible in one direction.
The principle was independently rediscovered by Thomas Edison on February 13, 1880. At the time Edison was carrying out research into why the filaments of his carbon-filament light bulbs nearly always burned out at the positive-connected end. He had a special bulb made with a metal plate sealed into the glass envelope, and he was able to confirm that an invisible current could be drawn from the glowing filament through the vacuum to the metal plate, but only when the plate was connected to the positive supply.
Edison devised a circuit where his modified light bulb more or less replaced the resistor in a DC voltmeter and on this basis was awarded a patent for it in 1883 (U.S. Patent 307,031). There was no apparent practical use for such device at the time, and the patent application was most likely simply a precaution in case someone else did find a use for the so-called “Edison Effect”.
About 20 years later, John Ambrose Fleming (scientific adviser to the Marconi Company and former Edison employee) realized that the Edison effect could be used as a precision radio detector. Fleming patented the first true thermionic diode in Britain [2] on November 16, 1904 (followed by U.S. Patent 803,684 in November 1905).
The principle of operation of crystal diodes was discovered in 1874 by the German scientist, Karl Ferdinand Braun.[3] Braun patented the crystal rectifier in 1899.[4] Braun’s discovery was further developed by Jagdish Chandra Bose into a useful device for radio detection.
The first actual radio receiver using a crystal diode was built by Greenleaf Whittier Pickard. Pickard received a patent for a silicon crystal detector on November 20, 1906[5] (U.S. Patent 836,531).
Other experimenters tried a variety of minerals and other substances, although by far the most popular was the Lead Sulfide mineral Galena. Although other substances offered slightly better performance, galena had the advantage of being cheap and easy to obtain, and was used almost exclusvely in home-built “crystal sets”, until the advent of inexpensive fixed germanium diodes in the 1950s.
At the time of their invention, such devices were known as rectifiers. In 1919, William Henry Eccles coined the term diode from Greek roots; dia means “through”, and ode (from ὅδος) means “path”.

Thermionic and gaseous stste diodes:-
Thermionic diodes are thermionic-valve devices (also known as vacuum tubes, tubes, or valves), which are arrangements of electrodes surrounded by a vacuum within a glass envelope. Early examples were fairly similar in appearance to incandescent light bulbs.
In thermionic valve diodes, a current through the heater filament indirectly heats the cathode, another internal electrode treated with a mixture of barium and strontium oxides, which are oxides of alkaline earth metals; these substances are chosen because they have a small work function. (Some valves use direct heating, in which a tungsten filament acts as both heater and cathode.) The heat causes thermionic emission of electrons into the vacuum. In forward operation, a surrounding metal electrode called the anode is positively charged so that it electrostatically attracts the emitted electrons. However, electrons are not easily released from the unheated anode surface when the voltage polarity is reversed. Hence, any reverse flow is negligible.
For much of the 20th century, thermionic valve diodes were used in analog signal applications, and as rectifiers in many power supplies. Today, valve diodes are only used in niche applications such as rectifiers in electric guitar and high-end audio amplifiers as well as specialized high-voltage equipment.

Semiconductor diodes:-

Most diodes today are based on semiconductor p-n junctions. In a p-n diode, conventional current is from the p-type side (the anode) to the n-type side (the cathode), but not in the opposite direction. Another type of semiconductor diode, the Schottky diode, is formed from the contact between a metal and a semiconductor rather than by a p-n junction.

Induction furnace

2009-08-24T23:48:00.000-07:00

An induction furnace is an electrical furnace in which the heat is applied by induction heating of a conductive medium (usually a metal) in a crucible placed in a water-cooled alternating current solenoid coil. The advantage of the induction furnace is a clean, energy-efficient and well-controllable melting process compared to most other means of metal melting. Most modern foundries use this type of furnace and now also more iron foundries are replacing cupolas with induction furnaces to melt cast iron, as the former emit lots of dust and other pollutants. Induction furnace capacities range from less than one kilogram to one hundred tonnes capacity, and are used to melt iron and steel, copper, aluminium, and precious metals. The one major drawback to induction furnace usage in a foundry is the lack of refining capacity; charge materials must be clean of oxidation products and of a known composition, and some alloying elements may be lost due to oxidation (and must be re-added to the melt).Operating frequencies range from utility frequency (50 or 60 Hz) to 400 kHz or higher, usually depending on the material being melted, the capacity(volume) of the furnace and the melting speed required. Generally the smaller the volume of the melts the higher the frequency of the furnace used; this is due to the skin depth which is a measure of the distance an alternating current can penetrate beneath the surface of a conductor. For the same conductivity the higher frequencies have a shallow skin depth - that is less penetration into the melt. Lower frequencies can generate stirring or turbulence in the metal.
A preheated 1-tonne furnace melting iron can melt cold charge to tapping readiness within an hour.
An operating induction furnace usually emits a hum or whine (due to magnetostriction), the pitch of which can be used by operators to identify whether the furnace is operating correctly, or at what power level.

Arc furnace

2009-08-24T23:36:00.000-07:00

An electric arc furnace (EAF) is a furnace that heats charged material by means of an electric arc.
Arc furnaces range in size from small units of approximately one ton capacity (used in foundries for producing cast iron products) up to about 400 ton units used for secondary steelmaking. Arc furnaces used in research laboratories and by dentists may have a capacity of only a few dozen grams. Electric arc furnace temperatures can be up to 1,800 degrees Celsius. Arc furnaces differ from induction furnaces in that the charge material is directly exposed to the electric arc, and the current in the furnace terminals passes through the charged material.
History:-
In the 19th century, a number of men had employed an electric arc to melt iron. Sir Humphry Davy conducted an experimental demonstration in 1810; welding was investigated by Pepys in 1815; Pinchon attempted to create an electrothermic furnace in 1853; and, in 1878 - 79, Sir William Siemens took out patents for electric furnaces of the arc type.
The first electric arc furnaces were developed by Paul Héroult, of France, with a commercial plant established in the United States in 1907. Initially "electric steel" was a specialty product for such uses as machine tools and spring steel. Arc furnaces were also used to prepare calcium carbide for use in carbide lamps. The Stessano electric furnace is an arc type furnace that usually rotates to mix the bath. The Girod furnace is similar to the Héroult furnace.
While EAFs were widely used in World War II for production of alloy steels, it was only later that electric steelmaking began to expand. The low capital cost for a mini-mill - around US$140-200 per ton of annual installed capacity, compared with US$1,000 per ton of annual installed capacity for an integrated steel mill - allowed mills to be quickly established in war-ravaged Europe, and also allowed them to successfully compete with the big United States steelmakers, such as Bethlehem Steel and U.S. Steel, for low-cost, carbon steel 'long products' (structural steel, rod and bar, wire and fasteners) in the U.S. market.
When Nucor - now one of the largest steel producers in the U.S.[1] - decided to enter the long products market in 1969, they chose to start up a mini-mill, with an EAF as its steelmaking furnace, soon followed by other manufacturers. Whilst Nucor expanded rapidly in the Eastern US, the companies that followed them into mini-mill operations concentrated on local markets for long products, where the use of an EAF allowed the plants to vary production according to local demand. This pattern was also followed globally, with EAF steel production primarily used for long products, while integrated mills, using blast furnaces and basic oxygen furnaces, cornered the markets for 'flat products' - sheet steel and heavier steel plate. In 1987, Nucor made the decision to expand into the flat products market, still using the EAF production method[2].
Construction:-
An electric arc furnace used for steelmaking consists of a refractory-lined vessel, usually water-cooled in larger sizes, covered with a retractable roof, and through which one or more graphite electrodes enter the furnace. The furnace is primarily split into three sections:
the shell, which consists of the sidewalls and lower steel 'bowl'; the hearth, which consists of the refractory that lines the lower bowl; the roof, which may be refractory-lined or water-cooled, and can be shaped as a section of a sphere, or as a frustum (conical section). The roof also supports the refractory delta in its centre, through which one or more graphite electrodes enter. The hearth may be hemispherical in shape, or in an eccentric bottom tapping furnace (see below), the hearth has the shape of a halved egg. In modern meltshops, the furnace is often raised off the ground floor, so that ladles and slag pots can easily be maneuvered under either end of the furnace. Separate from the furnace structure is the electrode support and electrical system, and the tilting platform on which the furnace rests. Two configurations are possible: the electrode supports and the roof tilt with the furnace, or are fixed to the raised platform.
A typical alternating current furnace has three electrodes. Electrodes are round in section, and typically in segments with threaded couplings, so that as the electrodes wear, new segments can be added. The arc forms between the charged material and the electrode, the charge is heated both by current passing through the charge and by the radiant energy evolved by the arc. The electrodes are automatically raised and lowered by a positioning system, which may use either electric winch hoists or hydraulic cylinders. The regulating system maintains approximately constant current and power input during the melting of the charge, even though scrap may move under the electrodes as it melts. The mast arms holding the electrodes carry heavy busbars, which may be hollow water-cooled copper pipes carrying current to the electrode holders. Modern systems use 'hot arms', where the whole arm carries the current, increasing efficiency. These can be made from copper-clad steel or aluminium. Since the electrodes move up and down automatically for regulation of the arc, and are raised to allow removal of the furnace roof, heavy water-cooled cables connect the bus tubes/arms with the transformer located adjacent to the furnace. To protect the transformer from heat, it is installed in a vault.
The furnace is built on a tilting platform so that the liquid steel can be poured into another vessel for transport. The operation of tilting the furnace to pour molten steel is called "tapping". Originally, all steelmaking furnaces had a tapping spout closed with refractory that washed out when the furnace was tilted, but often modern furnaces have an eccentric bottom tap-hole (EBT) to reduce inclusion of nitrogen and slag in the liquid steel. These furnaces have a taphole that passes vertically through the hearth and shell, and is set off-centre in the narrow 'nose' of the egg-shaped hearth. It is filled with refractory sand, such as olivine, when it is closed off. Modern plants may have two shells with a single set of electrodes that can be transferred between the two; one shell preheats scrap while the other shell is utilised for meltdown. Other DC-based furnaces have a similar arrangement, but have electrodes for each shell and one set of electronics.
AC furnaces usually exhibit a pattern of hot and cold-spots around the hearth perimeter, with the cold-spots located between the electrodes. Modern furnaces mount oxygen-fuel burners in the sidewall and use them to provide chemical energy to the cold-spots, making the heating of the steel more uniform. Additional chemical energy is provided by injecting oxygen and carbon into the furnace, historically this was done through lances in the slag door, noways this is mainly done through multiple wall-mounted injection units.
A mid-sized modern steelmaking furnace would have a transformer rated about 60,000,000 volt-amperes (60 MVA), with a secondary voltage between 400 and 900 volts and a secondary current in excess of 44,000 amperes. In a modern shop such a furnace would be expected to produce a quantity of 80 metric tonnes of liquid steel in approximately 60 minutes from charging with cold scrap to tapping the furnace. In comparison, basic oxygen furnaces can have a capacity of 150-300 tonnes per batch, or 'heat', and can produce a heat in 30-40 minutes. Enormous variations exist in furnace design details and operation, depending on the end product and local conditions, as well as ongoing research to improve furnace efficiency - the largest scrap-only furnace (in terms of tapping weight and transformer rating) is in Turkey, with a tap weight of 300 metric tonnes and a transformer of 300 MVA.
To produce a ton of steel in an electric arc furnace requires approximately 400 kilowatt-hours per short ton of electricity, or about 440kWh per metric tonne; the theoretical minimum amount of energy required to melt a tonne of scrap steel is 300kWh (melting point 1520°C/2768°F). Therefore, the 300-tonne, 300 MVA EAF mentioned above will require approximately 132 MWh of energy to melt the steel, and a 'power-on time' (the time that steel is being melted with an arc) of approximately 37 minutes, allowing for the power factor. Electric arc steelmaking is only economical where there is plentiful electricity, with a well-developed electrical grid.
Operation:-
Scrap metal is delivered to a scrap bay, located next to the melt shop. Scrap generally comes in two main grades: shred (whitegoods, cars and other objects made of similar light-gauge steel) and heavy melt (large slabs and beams), along with some direct reduced iron (DRI) or pig iron for chemical balance. Some furnaces melt almost 100% DRI. The scrap is loaded into large buckets called baskets, with 'clamshell' doors for a base. Care is taken to layer the scrap in the basket to ensure good furnace operation; heavy melt is placed on top of a light layer of protective shred, on top of which is placed more shred. These layers should be present in the furnace after charging. After loading, the basket may pass to a scrap pre-heater, which uses hot furnace off-gases to heat the scrap and recover energy, increasing plant efficiency. The scrap basket is then taken to the melt shop, the roof is swung off the furnace, and the furnace is charged with scrap from the basket. Charging is one of the more dangerous operations for the EAF operators. There is a lot of energy generated by multiple tonnes of falling metal; any liquid metal in the furnace is often displaced upwards and outwards by the solid scrap, and the grease and dust on the scrap is ignited if the furnace is hot, resulting in a fireball erupting. In some twin-shell furnaces, the scrap is charged into the second shell while the first is being melted down, and pre-heated with off-gas from the active shell. Other operations are continuous charging - pre-heating scrap on a conveyor belt, which then discharges the scrap into the furnace proper, or charging the scrap from a shaft set above the furnace, with off-gases directed through the shaft. Other furnaces can be charged with hot (molten) metal from other operations. After charging, the roof is swung back over the furnace and meltdown commences. The electrodes are lowered onto the scrap, an arc is struck and the electrodes are then set to bore into the layer of shred at the top of the furnace. Lower voltages are selected for this first part of the operation to protect the roof and walls from excessive heat and damage from the arcs. Once the electrodes have reached the heavy melt at the base of the furnace and the arcs are shielded by the scrap, the voltage can be increased and the electrodes raised slightly, lengthening the arcs and increasing power to the melt. This enables a molten pool to form more rapidly, reducing tap-to-tap times. Oxygen is also lanced into the scrap, combusting or cutting the steel, and extra chemical heat is provided by wall-mounted oxygen-fuel burners. Both processes accelerate scrap meltdown. An important part of steelmaking is the formation of slag, which floats on the surface of the molten steel. Slag usually consists of metal oxides, and acts as a destination for oxidised impurities, as a thermal blanket (stopping excessive heat loss) and helping to reduce erosion of the refractory lining. For a furnace with basic refractories, which includes most carbon steel-producing furnaces, the usual slag formers are calcium oxide (CaO, in the form of burnt lime) and magnesium oxide (MgO, in the form of dolomite and magnesite). These slag formers are either charged with the scrap, or blown into the furnace during meltdown. Another major component of EAF slag is iron oxide from steel combusting with the injected oxygen. Later in the heat, carbon (in the form of coke or coal) is lanced into this slag layer, reacting with the iron oxide to form metallic iron and carbon monoxide gas, which then causes the slag to foam, allowing greater thermal efficiency, and better arc stability and electrical efficiency. The slag blanket also covers the arcs, preventing damage to the furnace roof and sidewalls from radiant heat. Once flat bath conditions are reached, i.e. the scrap has been completely melted down, another bucket of scrap can be charged into the furnace and melted down, although EAF development is moving towards single-charge designs. After the second charge is completely melted, refining operations take place to check and correct the steel chemistry and superheat the melt above its freezing temperature in preparation for tapping. More slag formers are introduced and more oxygen is lanced into the bath, burning out impurities such as silicon, sulfur, phosphorus, aluminium, manganese and calcium and removing their oxides to the slag. Removal of carbon takes place after these elements have burnt out first, as they have a greater affinity for oxygen. Metals that have a poorer affinity for oxygen than iron, such as nickel and copper, cannot be removed through oxidation and must be controlled through scrap chemistry alone, such as introducing the direct reduced iron and pig iron mentioned earlier. A foaming slag is maintained throughout, and often overflows the furnace to pour out of the slag door into the slag pit. Temperature sampling and chemical sampling (in the form of a 'chill' - a small, solidified sample of the steel) take place via automatic lances. Once the temperature and chemistry are correct, the steel is tapped out into a preheated ladle through tilting the furnace. As soon as slag is detected during tapping the furnace is rapidly tilted back towards the deslagging side, minimising slag carryover into the ladle. During tapping some alloy additions are introduced into the metal stream. Often, a few tonnes of liquid steel and slag is left in the furnace in order to form a 'hot heel', which helps preheat the next charge of scrap and accelerate its meltdown. During and after tapping, the furnace is 'turned around': the slag door is cleaned of solidified slag, repairs may take place, and electrodes are inspected for damage or lengthened through the addition of new segments; the taphole is filled with sand at the completion of tapping. For a 90-tonne, medium-power furnace, the whole process will usually take about 60-70 minutes from the tapping of one heat to the tapping of the next (the tap-to-tap time).

Spot welding

2009-08-24T23:13:00.000-07:00

Spot welding is a process in which contacting metal surfaces are joined by the heat obtained from resistance to electric current flow. Work-pieces are held together under pressure exerted by electrodes. Typically the sheets are in the 0.5-3.0 mm thickness range. The process uses two shaped copper alloy electrodes to concentrate welding current into a small "spot" and to simultaneously clamp the sheets together. Forcing a large current through the spot will melt the metal and form the weld. The attractive feature of spot welding is a lot of energy can be delivered to the spot in a very short time (ten to one hundred milliseconds[citation needed]). That permits the welding to occur without excessive heating to the rest of the sheet.
The amount of heat (energy) delivered to the spot is determined by the resistance between the electrodes and the amplitude and duration of the current. The amount of energy is chosen to match the sheet's material properties, its thickness, and type of electrodes. Applying too little energy won't melt the metal or will make a poor weld. Applying too much energy will melt too much metal and make a hole rather than a weld. Another attractive feature of spot welding is the energy delivered to the spot can be controlled to produce reliable welds.
Projection welding is a modification of spot welding. In this process the weld is localized by means of raised sections, or projections, on one or both of the workpieces to be joined. heat is concentrated at the projections, which permits the welding of heavier sections or the closer spacing of welds. The projections can also serve as a means of positioning the workpieces. Projection welding is often used to weld studs, nuts, and other screw machine parts to metal plate. It's also frequently used to join crossed wires and bars. This is another high-production process, and multiple projection welds can be arranged by suitable
Applications:-
Spot welding is typically used when welding particular types of sheet metal. Thicker stock is more difficult to spot weld because the heat flows into the surrounding metal more easily. Spot welding can be easily identified on many sheet metal goods, such as metal buckets. Aluminum alloys can also be spot welded. However, their much higher thermal conductivity and electrical conductivity mean that up to three times higher welding currents are needed. This requires larger, more powerful, and more expensive welding transformers.
Perhaps the most common application of spot welding is in the automobile manufacturing industry, where it is used almost universally to weld the sheet metal to form a car. Spot welders can also be completely automated, and many of the industrial robots found on assembly lines are spot welders (the other major use for robots being painting).
Spot welding is also used is in the orthodontist's clinic, where small scale spot welding equipment is used when resizing metal "molar bands" used in orthodontics.
Another application is spot welding straps to nickel-cadmium or nickel-metal-hydride cells in order to make batteries. The cells are joined by spot welding thin nickel straps to the battery terminals. Spot welding can keep the battery from getting too hot, as might happen if conventional soldering were done.
Good design practice must always allow for adequate accessibility. Connecting surfaces should be free of contaminants, such as scale, oil, and dirt, for quality welds. Metal thickness is generally not a factor in determining good welds.
Processing and Equipment:-
Spot welding involves three stages; the first of which involves the electrodes being brought to the surface of the metal and applying a slight amount of pressure. The current from the electrodes is then applied briefly after which the current is removed but the electrodes remain in place in order for the material to cool. Weld times range from 0.01 sec to 0.63 sec depending on the thickness of the metal, the electrode force and the diameter of the electrodes themselves.
The equipment used in the spot welding process consists of tool holders and electrodes. The tool holders function as a mechanism to hold the electrodes firmly in place and also support optional water hoses which cool the electrodes during welding. Tool holding methods include a paddle-type, light duty, universal, and regular offset. The electrodes generally are made of a low resistance alloy, usually copper, and are designed in many different shapes and sizes depending on the application needed.
The two materials being welded together are known as the workpieces and must conduct electricity. The width of the workpieces is limited by the throat length of the welding apparatus and ranges typically from 5 to 50 inches. Workpiece thickness can range from 0.008in. to 1.25in.
After the current is removed from the workpiece, it is cooled via the coolant holes in the center of the electrodes. Both water and a brine solution may be used as coolants in spot welding mechanisms.
Tool styles:-
Electrodes used in spot welding can vary greatly with different applications. Each tool style has a different purpose. Radius style electrodes are used for high heat applications, electrodes with a truncated tip for high pressure, eccentric electrodes for welding corners, offset eccentric tips for reaching into corners and small spaces, and finally offset truncated for reaching into the workpiece itself.
Effects:-
The spot welding process tends to harden the material, cause it to warp, reduce the materials fatigue strength, and may stretch the material as well as anneal it. The physical effects of spot welding include internal cracking, surface cracks and a bad appearance. The chemical properties affected include the metal's internal resistance and its corrosive properties.
Electrical notes:-
The basic spot welder consists of a power supply, an energy storage unit (e.g., a capacitor bank), a switch, a welding transformer, and the welding electrodes. The energy storage element allows the welder to deliver high instantaneous power levels. If the power demands are not high, then the energy storage element isn't needed. The switch causes the stored energy to be dumped into the welding transformer. The welding transformer steps down the voltage and steps up the current. An important feature of the transformer is it reduces the current level that the switch must handle. The welding electrodes are part of the transformer's secondary circuit. There is also a control box that manages the switch and may monitor the welding electrode voltage or current.
The resistance presented to the welder is complicated. There is the resistance of secondary winding, the cables, and the welding electrodes. There is also the contact resistance between the welding electrodes and the workpiece. There is the resistance of the workpieces, and the contact resistance between the workpieces.
At the beginning of the weld, the contact resistances are usually high, so most of the initial energy will be dissipated there. That heat and the clamping force will soften and smooth out the material at the electrode-material interface and make better contact (that is, lower the contact resistance). Consequently, more electrical energy will go into the workpiece and the junction resistance of the two workpieces. As electrical energy is delivered to the weld and causes the temperature to rise, the electrodes and the workpiece are conducting that heat away. The goal is to apply enough energy so that a portion of material within the spot melts without having the entire spot melt. The perimeter of the spot will conduct away a lot of heat and keep the perimeter at a lower temperature. The interior of the spot has less heat conducted away, so it melts first. If the welding current is applied too long, the entire spot melts, the material runs out or otherwise fails, and the "weld" becomes a hole.
The voltage needed for welding depends on the resistance of the material to be welded, the sheet thickness and desired size of the nugget. When welding a common combination like 1.0 + 1.0 mm sheet steel, the voltage between the electrodes is only about 1.5 V at the start of the weld but can fall as low as 1 V at the end of the weld. This decrease in voltage results from the reduction in resistance caused by the workpiece melting. The open circuit voltage from the transformer is higher than this, typically in the 5-10 V range[citation needed], but there is a large voltage drop in the electrodes and secondary side of the transformer when the circuit is closed[citation needed].The resistance of the weld spot changes as it flows and liquefies. Modern welding equipment can monitor and adjust the weld in real-time to ensure a consistent weld. The equipment may seek to control different variables during the weld, such as current, voltage, power, or energy.
A spot welding machine normally consists of tool and electrodes, which are mechanisms for making and holding contact at the weld. Tool holders have two functions: to hold the electrode firmly in place and to support water hoses that provide cooling of the electrodes.
Weld currents can range from 4000 to 24,000 amps for different types of 1010 mild steel.
Safety Factors:-
Spot welding can be extremely dangerous with the use of large amounts of current and heat. Always be sure to clamp the electrodes tightly, protect your eyes from the intense light given off in the welding process and protect your hands because the workpiece can get very hot during welding. Take all safety precautions necessary before, during and after spot welding.

Air conditioner

2009-08-23T08:13:00.000-07:00

Air conditioning refers to the cooling and dehumidification of indoor air for thermal comfort. In a broader sense, the term can refer to any form of cooling, heating, ventilation or disinfection that modifies the condition of air. An air conditioner (often referred to as AC or air con.) is an appliance, system, or mechanism designed to stabilise the air temperature and humidity within an area (used for cooling as well as heating depending on the air properties at a given time), typically using a refrigeration cycle but sometimes using evaporation, commonly for comfort cooling in buildings and motor vehicles.
The concept of air conditioning is known to have been applied in Ancient Rome, where aqueduct water was circulated through the walls of certain houses to cool them. Similar techniques in medieval Persia involved the use of cisterns and wind towers to cool buildings during the hot season. Modern air conditioning emerged from advances in chemistry during the 19th century, and the first large-scale electrical air conditioning was invented and used in 1902 by Willis Haviland Carrier.
History:-

While moving heat via machinery to provide air conditioning is a relatively modern invention, the cooling of buildings is not. Wealthy ancient Romans circulated aqueduct water through walls to cool their luxurious houses.[citation needed]
The 2nd century Chinese inventor Ding Huan (fl. 180) of the Han Dynasty invented a rotary fan for air conditioning, with seven wheels 3 m (10 ft) in diameter and manually powered.In 747, Emperor Xuanzong (r. 712–762) of the Tang Dynasty (618–907) had the Cool Hall (Liang Tian) built in the imperial palace, which the Tang Yulin describes as having water-powered fan wheels for air conditioning as well as rising jet streams of water from fountains.[3] During the subsequent Song Dynasty (960–1279), written sources mentioned the air conditioning rotary fan as even more widely used.
Medieval Persia had buildings that used cisterns and wind towers to cool buildings during the hot season: cisterns (large open pools in central courtyards, not underground tanks) collected rain water; wind towers had windows that could catch wind and internal vanes to direct the airflow down into the building, usually over the cistern and out through a downwind cooling tower. Cistern water evaporated, cooling the air in the building.
Ventilators were invented in medieval Egypt and were widely used in many houses throughout Cairo during the Middle Ages. These ventilators were later described in detail by Abd al-Latif al-Baghdadi in 1200, who reported that almost every house in Cairo has a ventilator, and that they cost anywhere from 1 to 500 dinars depending on their sizes and shapes. Most ventilators in the city were oriented towards the Qibla, as was the city in general.
In the 1600s Cornelius Drebbel demonstrated "turning Summer into Winter" for James I of England by adg salt to water.
In 1758, Benjamin Franklin and John Hadley, professor of chemistry at Cambridge University, conducted an experiment to explore the principle of evaporation as a means to rapidly cool an object. Franklin and Hadley confirmed that evaporation of highly volatile liquids such as alcohol and ether, could be used to drive down the temperature of an object past the freezing point of water. They conducted their experiment with the bulb of a mercury thermometer as their object and with a bellows used to "quicken" the evaporation; they lowered the temperature of the thermometer bulb down to 7°F while the ambient temperature was 65°F. Franklin noted that soon after they passed the freezing point of water (32°F) a thin film of ice formed on the surface of the thermometer's bulb and that the ice mass was about a quarter inch thick when they stopped the experiment upon reaching 7°F. Franklin concluded, "From this experiment, one may see the possibility of freezing a man to death on a warm summer's day".
In 1820, British scientist and inventor Michael Faraday discovered that compressing and liquefying ammonia could chill air when the liquefied ammonia was allowed to evaporate. In 1842, Florida physician John Gorrie used compressor technology to create ice, which he used to cool air for his patients in his hospital in Apalachicola, Florida.He hoped eventually to use his ice-making machine to regulate the temperature of buildings. He even envisioned centralized air conditioning that could cool entire cities. Though his prototype leaked and performed irregularly, Gorrie was granted a patent in 1851 for his ice-making machine. His hopes for its success vanished soon afterwards when his chief financial backer died; Gorrie did not get the money he needed to develop the machine. According to his biographer, Vivian M. Sherlock, he blamed the "Ice King", Frederic Tudor, for his failure, suspecting that Tudor had launched a smear campaign against his invention. Dr. Gorrie died impoverished in 1855 and the idea of air conditioning faded away for 50 years.
Early commercial applications of air conditioning were manufactured to cool air for industrial processing rather than personal comfort. In 1902 the first modern electrical air conditioning was invented by Willis Haviland Carrier in Syracuse, NY. Designed to improve manufacturing process control in a printing plant, his invention controlled not only temperature but also humidity. The low heat and humidity were to help maintain consistent paper dimensions and ink alignment. Later Carrier's technology was applied to increase productivity in the workplace, and The Carrier Air Conditioning Company of America was formed to meet rising demand. Over time air conditioning came to be used to improve comfort in homes and automobiles. Residential sales expanded dramatically in the 1950s.
In 1906, Stuart W. Cramer of Charlotte, North Carolina, USA, was exploring ways to add moisture to the air in his textile mill. Cramer coined the term "air conditioning", using it in a patent claim he filed that year as an analogue to "water conditioning", then a well-known process for making textiles easier to process. He combined moisture with ventilation to "condition" and change the air in the factories, controlling the humidity so necessary in textile plants. Willis Carrier adopted the term and incorporated it into the name of his company. This evaporation of water in air, to provide a cooling effect, is now known as evaporative cooling.
The first air conditioners and refrigerators employed toxic or flammable gases like ammonia, methyl chloride, and propane which could result in fatal accidents when they leaked. Thomas Midgley, Jr. created the first chlorofluorocarbon gas, Freon, in 1928.
Freon is a trademark name of DuPont for any Chlorofluorocarbon (CFC), Hydrogenated CFC (HCFC), or Hydrofluorocarbon (HFC) refrigerant, the name of each including a number indicating molecular composition (R-11, R-12, R-22, R-134A). The blend most used in direct-expansion home and building comfort cooling is an HCFC known as R-22. It is to be phased out for use in new equipment by 2010 and completely discontinued by 2020. R-12 was the most common blend used in automobiles in the US until 1994 when most changed to R-134A. R-11 and R-12 are no longer manufactured in the US for this type of application, the only source for air conditioning purchase being the cleaned and purified gas recovered from other air conditioner systems. Several non-ozone depleting refrigerants have been developed as alternatives, including R-410A, invented by Honeywell (formerly AlliedSignal) in Buffalo NY and sold under the Genetron (R) AZ-20 name it was first commercially used by Carrier under the brand name Puron.
Innovation in air conditioning technologies continue, with much recent emphasis placed on energy efficiency, and improving indoor air quality. Reducing climate change impacts is an important area of innovation, because in addition to greenhouse gas emissions associated with energy use, CFCs, HCFCs and HFCs are potent greenhouse gases when leaked to the atmosphere. For example, R-22 (also known as HCFC-22) has a global warming potential about 1,800 times higher than CO2. As an alternative to conventional refrigerants, natural alternatives like CO2 (R-744) have been proposed.
Air conditioning applications:-
Air conditioning engineers broadly divide air conditioning applications into comfort and process.
Comfort applications aim to provide a building indoor environment that remains relatively constant in a range preferred by humans despite changes in external weather conditions or in internal heat loads.
Air conditioning makes deep plan buildings feasible, for otherwise they'd have to be built narrower or with light wells so that inner spaces receive sufficient outdoor air via natural ventilation. Air conditioning also allows buildings to be taller since wind speed increases significantly with altitude making natural ventilation impractical for very tall buildings. Comfort applications for various building types are quite different and may be categorized as
Low-Rise Residential buildings, including single family houses, duplexes, and small apartment buildings High-Rise Residential buildings, such as tall dormitories and apartment blocks Commercial buildings, which are built for commerce, including offices, malls, shopping centers, restaurants, etc. Institutional buildings, which includes hospitals, governmental, academic, and so on. Industrial spaces where thermal comfort of workers is desired. In addition to buildings, air conditioning can be used for many types of transportation - motor-cars and other land vehicles, trains, ships, aircraft, and spacecraft.
Process applications aim to provide a suitable environment for a process being carried out, regardless of internal heat and humidity loads and external weather conditions. Although often in the comfort range, it is the needs of the process that determine conditions, not human preference. Process applications include these:
Hospital operating theatres, in which air is filtered to high levels to reduce infection risk and the humidity controlled to limit patient dehydration. Although temperatures are often in the comfort range, some specialist procedures such as open heart surgery require low temperatures (about 18 °C, 64 °F) and others such as neonatal relatively high temperatures (about 28 °C, 82 °F). Cleanrooms for the production of integrated circuits, pharmaceuticals, and the like, in which very high levels of air cleanliness and control of temperature and humidity are required for the success of the process. Facilities for breeding laboratory animals. Since many animals normally only reproduce in spring, holding them in rooms at which conditions mirror spring all year can cause them to reproduce year-round. Aircraft air conditioning. Although nominally aimed at providing comfort for passengers and cooling of equipment, aircraft air conditioning presents a special challenge because of the changing density altitude associated with changes in altitude, humidity and temperature of the outside air[vague]. Data centers Textile factories Physical testing facilities Plants and farm growing areas Nuclear facilities Chemical and biological laboratories Mines Industrial environments Food cooking and processing areas In both comfort and process applications the objective may be to not only control temperature, but also humidity, air quality and air movement from space to space.
Humidity control:-
Refrigeration air conditioning equipment usually reduces the humidity of the air processed by the system. The relatively cold (below the dewpoint) evaporator coil condenses water vapor from the processed air, (much like an ice-cold drink will condense water on the outside of a glass), sending the water to a drain and removing water vapor from the cooled space and lowering the relative humidity. Since humans perspire to provide natural cooling by the evaporation of perspiration from the skin, drier air (up to a point) improves the comfort provided. The comfort air conditioner is designed to create a 40% to 60% relative humidity in the occupied space. In food retailing establishments large open chiller cabinets act as highly effective air dehumidifying units.
A specific type of air conditioner that is used only for dehumidifying is called a dehumidifier. A dehumidifier is different from a regular air conditioner in that both the evaporator and condensor coils are placed in the same air path, and the entire unit is placed in the environment that is intended to be conditioned (in this case dehumidified), rather than requiring the condensor coil to be outdoors. Having the condensor coil in the same air path as the evaporator coil produces warm, dehumidified air. The evaporator (cold) coil is placed first in the air path, dehumidifying the air exactly as a regular air conditioner does. The air next passes over the condensor coil re-warming the now dehumidified air. Note that the terms "condensor coil" and "evaporator coil" do not refer to the behavior of water in the air as it passes over each coil; instead they refer to the phases of the refrigeration cycle. Having the condensor coil in the main air path rather than in a separate, outdoor air path (as in a regular air conditioner) results in two consequences—the output air is warm rather than cold, and the unit is able to be placed anywhere in the environment to be conditioned, without a need to have the condensor outdoors.
Unlike a regular air conditioner, a dehumidifier will actually heat a room just as an electric heater that draws the same amount of power (watts) as the dehumidifier. A regular air conditioner transfers energy out of the room by means of the condensor coil, which is outside the room (outdoors). This is a thermodynamic system where the room serves as the system and energy is transferred out of the system. Conversely with a dehumidifier, no energy is transferred out of the thermodynamic system (room) because the air conditioning unit (dehumidifier) is entirely inside the room. Therefore all of the power consumed by the dehumidifier is energy that is input into the thermodynamic system (the room), and remains in the room (as heat). In addition, if the condensed water has been removed from the room, the amount of heat needed to boil that water has been added to the room. This is the inverse of adding water to the room with an evaporative cooler.
Dehumidifiers are commonly used in cold, damp climates to prevent mold growth indoors, especially in basements. They are also sometimes used in hot, humid climates for comfort because they reduce the humidity which causes discomfort (just as a regular air conditioner, but without cooling the room).
The engineering of physical and thermodynamic properties of gas-vapor mixtures is named Psychrometrics.
Healyh implications:-

A poorly maintained air-conditioning system can occasionally promote the growth and spread of microorganisms, such as Legionella pneumophila, the infectious agent responsible for Legionnaires' disease, or thermophilic actinomycetes,but as long as the air conditioner is kept clean these health hazards can be avoided. Conversely, air conditioning, including filtration, humidification, cooling, disinfection, etc., can be used to provide a clean, safe, hypoallergenic atmosphere in hospital operating rooms and other environments where an appropriate atmosphere is critical to patient safety and well-being. Air conditioning can have a positive effect on sufferers of allergies and asthma
In serious heat waves, air conditioning can save the lives of the elderly. Some local authorities even set up public cooling centers for the benefit of those without air conditioning at home.
Energy use:-

It should be noted that in a thermodynamically closed system, any energy input into the system that is being maintained at a set temperature (which is a standard mode of operation for modern air conditioners) requires that the energy removal rate from the air conditioner increases. This increase has the effect that for each unit of energy input into the system (say to power a light bulb in the closed system) this requires the air conditioner to remove that energy.In order to do that the air conditioner must increase its consumption by the inverse of its efficiency times the input of energy. As an example, presume that inside the closed system a 100 watt light bulb is activated, and the air conditioner has an efficiency of 200%. The air conditioner's energy consumption will increase by 50 watts to compensate for this, thus making the 100 W light bulb use a total of 150 W of energy.
It is typical for air conditioners to operate at "efficiencies" of significantly greater than 100%. [16].However it may be noted that the input (electrical) energy is of higher thermodynamic quality than the output which is basically thermal energy (heat dissipated), See Coefficient of performance.

Automobile air conditions:-
Air conditioner systems are designed to allow the driver and or passengers to feel more comfortable during uncomfortably warm humid or hot trips in a vehicle. Cars in hot climates often are fitted with air conditioning. There has been much debate and discussion on what the usage of an air conditioner does to the fuel efficiency of a vehicle. Factors such as wind resistance aerodynamics and engine power and weight have to be factored into finding the true variance between using the air conditioning system and not using it when figuring out difference in actual gas mileage. Other factors on the impact on the engine and an overall engine heat increase can have an impact on the cooling system of the vehicle.
The Packard Motor Car Company was the first automobile manufacturer to build air conditioners into its cars, beginning in 1939. These air conditioners were originally optional, and could be installed for an extra $274 (about $4,050 in 2007 dollars),[18] though they took up the entire trunk space and were not very efficient.

Gas-discharge lamps

2009-08-23T07:43:00.000-07:00

Gas-discharge lamps are a family of artificial light sources that generate light by sending an electrical discharge through an ionized gas, i.e. a plasma. The character of the gas discharge critically depends on the frequency or modulation of the current: see the entry on a frequency classification of plasmas. Typically, such lamps use a noble gas (argon, neon, krypton and xenon) or a mixture of these gases. Most lamps are filled with additional materials, like mercury, sodium, and/or metal halides. In operation the gas is ionized, and free electrons, accelerated by the electrical field in the tube, collide with gas and metal atoms. Some electrons circling around the gas and metal atoms are excited by these collisions, bringing them to a higher energy state. When the electron falls back to its original state, it emits a photon, resulting in visible light or ultraviolet radiation. Ultraviolet radiation is converted to visible light by a fluorescent coating on the inside of the lamp's glass surface for some lamp types. The fluorescent lamp is perhaps the best known gas-discharge lamp.
Gas-discharge lamps offer long life and high light efficiency, but are more complicated to manufacture, and they require electronics to provide the correct current flow through the gas.
History:-
Francis Hauksbee first demonstrated a gas-discharge lamp in 1705. He showed that an evacuated or partially evacuated glass globe, while charged by static electricity could produce a light bright enough to read by. The phenomenon of electric arc was first described by Vasily V. Petrov, a Russian scientist, in 1802; sir Humphry Davy demonstrated in the same year the electric arc at the Royal Institution of Great Britain. Since then, discharge light sources have been researched because they create light from electricity considerably more efficiently than incandescent light bulbs.
Later it was discovered that the arc discharge could be optimized by using an inert gas instead of air as a medium. Therefore noble gases neon, argon, krypton or xenon were used, as well as carbon dioxide historically.
The introduction of the metal vapor lamp, including various metals within the discharge tube, was a later advance. The heat of the gas discharge vaporized some of the metal and the discharge is then produced almost exclusively by the metal vapor. The usual metals are sodium and mercury owing to their high vapor pressures that increase efficiency of visible spectrum emission.
One hundred years of research later led to lamps without electrodes which are instead energized by microwave or radio frequency sources. In addition, light sources of much lower output have been created, extending the applications of discharge lighting to home or indoor use.
Color:-
Each gas, depending on its atomic structure emits certain wavelengths which translates in different colors of the lamp. As a way of evaluating the ability of a light source to reproduce the colors of various objects being lit by the source, the International Commission on Illumination (CIE) introduced the color rendering index. Some gas-discharge lamps have a relatively low CRI, which means colors they illuminate appear substantially different than they do under sunlight or other high-CRI illumination.
Low pressure discharge lamps:-
Fluorescent lamps, the most common lamp in office lighting and many other applications, produces up to 100 lumens/watt Low pressure sodium lamps, the most efficient gas-discharge lamp type, producing up to 200 lumens/watt, but at the expense of very poor color rendering. The almost monochromatic yellow light is only acceptable for street lighting and similar applications.
High pressure discharge lamps:-

Metal halide lamps. These lamps produce almost white light, and attain 100 lumen/watt light output. Applications include indoor lighting of high buildings, parking lots, shops, sport terrains. High pressure sodium lamps, producing up to 150 lumens/watt. These lamps produce a broader light spectrum than the low pressure sodium lamps. Also used for street lighting, and for artificial photoassimilation for growing plants High pressure mercury-vapor lamps. This lamp type is the oldest high pressure lamp type, being replaced in most applications by the metal halide lamp and the high pressure sodium lamp.
High-intensity discharge lamps:-

A high-intensity discharge (HID) lamp is a type of electrical lamp which produces light by means of an electric arc between tungsten electrodes housed inside a translucent or transparent fused quartz or fused alumina arc tube. This tube is filled with both gas and metal salts. The gas facilitates the arc's initial strike. Once the arc is started, it heats and evaporates the metal salts forming a plasma, which greatly increases the intensity of light produced by the arc and reduces its power consumption. High-intensity discharge lamps are a type of arc lamp.
Compared with fluorescent and incandescent lamps, HID lamps have higher luminous efficacy since a greater proportion of their radiation is in visible light as opposed to heat. Their overall luminous efficacy is also much higher: they give a greater amount of light output per watt of electricity input.
Construction:-
Various different types of chemistry are used in the arc tubes of HID lamps, depending on the desired characteristics of light intensity, correlated color temperature, color rendering index (CRI), energy efficiency, and lifespan. Varieties of HID lamp include:
Mercury vapor lamps Metal halide (MH) lamps Ceramic MH lamps Sodium vapor lamps Xenon short-arc lamps. Ultra-High Performance (UHP) The light-producing element of these lamp types is a well-stabilized arc discharge contained within a refractory envelope arc tube with wall loading in excess of 3 W/cm² (19.4 W/in²).
Mercury vapor lamps were the first commercially available HID lamps. Originally they produced a bluish-green light, but more recent versions can produce light with a less pronounced color tint. However, mercury vapor lamps are falling out of favor and being replaced by sodium vapor and metal halide lamps.
Metal halide and ceramic metal halide lamps can be made to give off neutral white light useful for applications where normal color appearance is critical, such as TV and movie production, indoor or nighttime sports games, automotive headlamps, and aquarium lighting.
Low-pressure sodium vapor lamps are extremely efficient. They produce a deep yellow-orange light and have an effective CRI of nearly zero; items viewed under their light appear monochromatic. This makes them particularly effective as photographic safe lights. High-pressure sodium lamps tend to produce a much whiter light, but still with a characteristic orange-pink cast. New color-corrected versions producing a whiter light are now available, but some efficiency is sacrificed for the improved color.
Like fluorescent lamps, HID lamps require a ballast to start and maintain their arcs. The method used to initially strike the arc varies: mercury vapor lamps and some metal halide lamps are usually started using a third electrode near one of the main electrodes while other lamp styles are usually started using pulses of high voltage.
Applications:-
HID lamps are typically used when high levels of light over large areas are required, and when energy efficiency and/or light intensity are desired. These areas include gymnasiums, large public areas, warehouses, movie theaters, football stadiums[1], outdoor activity areas, roadways, parking lots, and pathways. More recently, HID lamps, especially metal halide, have been used in small retail and residential environments. HID lamps have made indoor gardening practical, particularly for plants that require a good deal of high intensity sunlight; HID lamps are a common choice of light source for marijuana growing operations. They are also used to reproduce tropical intensity sunlight for indoor aquariums. Ultra-High Performance (UHP) HID lamps are used in LCD or DLP projection TV sets or projection displays.
Most HID lamps produce significant UV radiation, and require UV-blocking filters to prevent UV-induced degradation of lamp fixture components and fading of dyed items illuminated by the lamp. Exposure to HID lamps operating with faulty or absent UV-blocking filters causes injury to humans and animals, such as sunburn and arc eye. Many HID lamps are designed so as to quickly extinguish if their outer UV-shielding glass envelope is broken.
Beginning in the early 1990s, HID lamps have been employed in motor vehicle headlamps. This application has met with mixed responses from motorists, who appreciate the improved nighttime visibility from HID headlamps but object to the glare they can cause. Internationalized European vehicle regulations require such headlamps to be equipped with lens cleaners and an automatic self-leveling system to keep the beams aimed correctly regardless of vehicle load and attitude, but no such devices are required in North America, where inherently more glaring beam patterns are also permitted. Retrofitting HID bulbs in headlamps not originally designed to accept them results in extremely high levels of glare, and is illegal throughout most of the world.
HID lamps are used in high-performance bicycle headlamps as well as flashlights and other portable lights, because they produce a great amount of light per unit of power. As the HID lights use less than half the power of an equivalent tungsten-halogen light, a significantly smaller and lighter-weight power supply can be used.
HID lamps are also used on many general aviation aircraft for landing and taxi lights.
End of life:-
Factors of wear come mostly from on/off cycles versus the total on time. The highest wear occurs when the HID burner is ignited while still hot and before the metallic salts have recrystallized.
At the end of life, many types of high-intensity discharge lamps exhibit a phenomenon known as cycling. These lamps can be started at a relatively low voltage. As they heat up during operation, however, the internal gas pressure within the arc tube rises and a higher voltage is required to maintain the arc discharge. As a lamp gets older, the voltage necessary to maintain the arc eventually rises to exceed the voltage provided by the electrical ballast. As the lamp heats to this point, the arc fails and the lamp goes out. Eventually, with the arc extinguished, the lamp cools down again, the gas pressure in the arc tube is reduced, and the ballast can once again cause the arc to strike. The effect of this is that the lamp glows for a while and then goes out, repeatedly.
More sophisticated ballast designs detect cycling and give up attempting to start the lamp after a few cycles. If power is removed and reapplied, the ballast will make a new series of startup attempts.
Sometimes the quartz tube containing mercury can explode in UHP lamps, especially when it is defective or weakened by many on/off cycles, or when pressure is excessive due to high temperature. When that happens, up to 30 mg vaporized mercury is released into atmosphere. It can be potentially toxic when indoors. A typical scenario is a failure of UHP HID lamp in front of rear LCD projection TV sets or computer displays. Some vendors recommend use of a mercury vacuum cleaner or respirator when dealing with bulb rupture due to risks of mercury vapors. They also require a special waste disposal.

Capacitor

2009-08-23T07:36:00.000-07:00

A capacitor or condenser is a passive electronic component consisting of a pair of conductors separated by a dielectric. When a voltage potential difference exists between the conductors, an electric field is present in the dielectric. This field stores energy and produces a mechanical force between the plates. The effect is greatest between wide, flat, parallel, narrowly separated conductors.
An ideal capacitor is characterized by a single constant value, capacitance, which is measured in farads. This is the ratio of the electric charge on each conductor to the potential difference between them. In practice, the dielectric between the plates passes a small amount of leakage current. The conductors and leads introduce an equivalent series resistance and the dielectric has an electric field strength limit resulting in a breakdown voltage.
Capacitors are widely used in electronic circuits to block the flow of direct current while allowing alternating current to pass, to filter out interference, to smooth the output of power supplies, and for many other purposes. They are used in resonant circuits in radio frequency equipment to select particular frequencies from a signal with many frequencies.
History:-
In October 1745, Ewald Georg von Kleist of Pomerania in Germany found that charge could be stored by connecting a high voltage electrostatic generator by a wire to a volume of water in a hand-held glass jar.[1] Von Kleist's hand and the water acted as conductors and the jar as a dielectric (although details of the mechanism were incorrectly identified at the time). Von Kleist found that after removing the generator, touching the wire resulted in a painful spark. In a letter describing the experiment, he said "I would not take a second shock for the kingdom of France."[2] The following year, the Dutch physicist Pieter van Musschenbroek invented a similar capacitor, which was named the Leyden jar, after the University of Leyden where he worked.[3] Daniel Gralath was the first to combine several jars in parallel into a "battery" to increase the charge storage capacity.[citation needed]
Benjamin Franklin investigated the Leyden jar and proved that the charge was stored on the glass, not in the water as others had assumed.[citation needed] He also created the term "battery",[4][5] (as in a battery of cannon), subsequently applied to clusters of electrochemical cells.[6] Leyden jars were later to be made by coating the inside and outside of jars with metal foil, leaving a space at the mouth to prevent arcing between the foils.[citation needed] The earliest unit of capacitance was the 'jar', equivalent to about 1 nanofarad.[citation needed]
Leyden jars or more powerful devices employing flat glass plates alternating with foil conductors were used exclusively up until about 1900, when the invention of wireless (radio) created a demand for standard capacitors, and the steady move to higher frequencies required capacitors with lower inductance.[citation needed] A more compact construction began to be used of a flexible dielectric sheet such as oiled paper sandwiched between sheets of metal foil, rolled or folded into a small package.[citation needed]
Early capacitors were also known as condensers, a term that is still occasionally used today. The term was first used for this purpose by Alessandro Volta in 1782, with reference to the device's ability to store a higher density of electric charge than a normal isolated conductor.[citation needed]

Fuse

2009-08-23T02:39:00.000-07:00

A fuse consists of a metal strip or wire fuse element, of small cross-section compared to the circuit conductors, mounted between a pair of electrical terminals, and (usually) enclosed by a non-conducting and non-combustible housing. The fuse is arranged in series to carry all the current passing through the protected circuit. The resistance of the element generates heat due to the current flow. The size and construction of the element is (empirically) determined so that the heat produced for a normal current does not cause the element to attain a high temperature. If too high a current flows, the element rises to a higher temperature and either directly melts, or else melts a soldered joint within the fuse, opening the circuit.
When the metal conductor parts, an electric arc forms between the un-melted ends of the element. The arc grows in length until the voltage required to sustain the arc is higher than the available voltage in the circuit, terminating current flow. In alternating current circuits the current naturally reverses direction on each cycle, greatly enhancing the speed of fuse interruption. In the case of a current-limiting fuse, the arc voltage builds up quickly enough to essentially stop the fault current before the first peak of the ac waveform. This effect significantly limits damage to downstream protected devices.
The fuse element is made of zinc, copper, silver, aluminum, or alloys to provide stable and predictable characteristics. The fuse ideally would carry its rated current indefinitely, and melt quickly on a small excess. The element must not be damaged by minor harmless surges of current, and must not oxidize or change its behavior after possibly years of service.
The fuse elements may be shaped to increase heating effect. In large fuses, current may be divided between multiple strips of metal. A dual-element fuse may contain a metal strip that melts instantly on a short-circuit, and also contain a low-melting solder joint that responds to long-term overload of low values compared to a short-circuit. Fuse elements may be supported by steel or nichrome wires, so that no strain is placed on the element, but a spring may be included to increase the speed of parting of the element fragments.
The fuse element may be surrounded by air, or by materials intended to speed the quenching of the arc. Silica sand or non-conducting liquids may be used.

Insulator

2009-08-23T02:24:00.000-07:00

An insulator, also called a dielectric, is a material that resists the flow of electric current. An insulating material has atoms with tightly bonded valence electrons. These materials are used in parts of electrical equipment, also called insulators or insulation, intended to support or separate electrical conductors without passing current through themselves. The term is also used more specifically to refer to insulating supports that attach electric power transmission wires to utility poles or pylons.
Some materials such as glass or Teflon are very good electrical insulators. A much larger class of materials, for example rubber-like polymers and most plastics are still "good enough" to insulate electrical wiring and cables even though they may have lower bulk resistivity. These materials can serve as practical and safe insulators for low to moderate voltages (hundreds, or even thousands, of volts).
physics of conduction in solids:-

Electrical insulation is the absence of electrical conduction. Electronic band theory (a branch of physics) predicts that a charge will flow whenever there are states available into which the electrons in a material can be excited. This allows them to gain energy and thereby move through the conductor (usually a metal). If no such states are available, the material is an insulator.
Most (though not all, see Mott insulator) insulators are characterized by having a large band gap. This occurs because the "valence" band containing the highest energy electrons is full, and a large energy gap separates this band from the next band above it. There is always some voltage (called the breakdown voltage) that will give the electrons enough energy to be excited into this band. Once this voltage is exceeded, the material ceases being an insulator, and charge will begin to pass through it. However, it is usually accompanied by physical or chemical changes that permanently degrade the material's insulating properties.
Materials that lack electron conduction are insulators if they lack other mobile charges as well. For example, if a liquid or gas contains ions, then the ions can be made to flow as an electric current, and the material is a conductor. Electrolytes and plasmas contain ions and will act as conductors whether or not electron flow is involved.
Breakdown:-
Insulators suffer from the phenomenon of electrical breakdown. When any voltage applied across a length of insulating substance exceeds a threshold breakdown field for that substance, which equals the band gap energy, the insulator suddenly turns into a resistor, sometimes with catastrophic results. During electrical breakdown, any free charge carrier being accelerated by the strong e-field will have enough velocity to knock electrons from (ionize) any atom it strikes. These freed electrons and ions are in turn accelerated and strike other atoms, creating more charge carriers, in a chain reaction. Rapidly the insulator becomes filled with mobile carriers, and its resistance drops to a low level. In air, the outbreak of conductivity is called "corona discharge" or a "spark." Similar breakdown can occur within any insulator, even within the bulk solid of a material. Even a vacuum can suffer a sort of break down, but in this case the breakdown or vacuum arc involves charges ejected from the surface of metal electrodes rather than produced by the vacuum itself.
Uses:-

Insulators are commonly used as a flexible coating on electric wire and cable. Since air is an insulator, no other substance is needed to "keep the electricity within the wires." However, wires which touch each other will produce cross connections, short circuits, and fire hazards. In coaxial cable the center conductor must be supported exactly in the middle of the hollow shield in order to prevent EM wave reflections. And any wires which present voltages higher than 60V can cause human shock and electrocution hazards. Insulating coatings prevent all of these problems.
In electronic systems, printed circuit boards are made from epoxy plastic and fiberglass. The nonconductive boards support layers of copper foil conductors. In electronic devices, the tiny and delicate active components are embedded within nonconductive epoxy or phenolic plastics, or within baked glass or ceramic coatings.
In microelectronic components such as transistors and ICs, the silicon material is normally a conductor because of doping, but it can easily be selectively transformed into a good insulator by the application of heat and oxygen. Oxidized silicon is quartz, i.e. silicon dioxide.
In high voltage systems containing transformers and capacitors, liquid insulator oil is the typical method used for preventing sparks. The oil replaces the air in any spaces which must support significant voltage without electrical breakdown. Other methods of insulating high voltage systems are ceramic or glass wire holders and simply placing the wires with a large separation, using the air as insulation.
Material:-Insulators used for high-voltage power transmission are made from glass, porcelain, or composite polymer materials. Porcelain insulators are made from clay, quartz or alumina and feldspar, and are covered with a smooth glaze to shed water. Insulators made from porcelain rich in alumina are used where high mechanical strength is a criterion. Porcelain has a dielectric strength of about 4–10 kV/mm.Glass has a higher dielectric strength, but it attracts condensation and the thick irregular shapes needed for insulators are difficult to cast without internal strains.Some insulator manufacturers stopped making glass insulators in the late 1960s, switching to ceramic materials.
Recently, some electric utilities have begun converting to polymer composite materials for some types of insulators. These are typically composed of a central rod made of fibre reinforced plastic and an outer weathershed made of silicone rubber or EPDM. Composite insulators are less costly, lighter in weight, and have excellent hydrophobic capability. This combination makes them ideal for service in polluted areas. However, these materials do not yet have the long-term proven service life of glass and porcelain.
Design:-

The electrical breakdown of an insulator due to excessive voltage can occur in one of two ways:
Puncture voltage is the voltage across the insulator (when installed in its normal manner) which causes a breakdown and conduction through the interior of the insulator. The heat resulting from the puncture arc usually damages the insulator irreparably. Flashover voltage is the voltage which causes the air around or along the surface of the insulator to break down and conduct, causing a 'flashover' arc along the outside of the insulator. They are usually designed to withstand this without damage. Most high voltage insulators are designed with a lower flashover voltage than puncture voltage, so they will flashover before they puncture, to avoid damage.
Dirt, pollution, salt, and particularly water on the surface of a high voltage insulator can create a conductive path across it, causing leakage currents and flashovers. The flashover voltage can be more than 50% lower when the insulator is wet. High voltage insulators for outdoor use are shaped to maximize the length of the leakage path along the surface from one end to the other, called the creepage length, to minimize these leakage currents.[3] To accomplish this the surface is molded into a series of corrugations or concentric disk shapes. These usually include one or more sheds; downward facing cup-shaped surfaces that act as umbrellas to ensure that the part of the surface leakage path under the 'cup' stays dry in wet weather. Minimum creepage distances are 20–25 mm/kV, but must be increased in high pollution or airborne sea-salt areas.

Earthing

2009-08-23T02:02:00.000-07:00

In electricity supply systems, an earthing system defines the electrical potential of the conductors relative to that of the Earth's conductive surface. The choice of earthing system has implications for the safety and electromagnetic compatibility of the power supply. Note that regulations for earthing (grounding) systems vary considerably between different countries.
A protective earth (PE) connection ensures that all exposed conductive surfaces are at the same electrical potential as the surface of the Earth, to avoid the risk of electrical shock if a person touches a device in which an insulation fault has occurred. It ensures that in the case of an insulation fault (a "short circuit"), a very high current flows, which will trigger an overcurrent protection device (fuse, circuit breaker) that disconnects the power supply.
A functional earth connection serves a purpose other than providing protection against electrical shock. In contrast to a protective earth connection, a functional earth connection may carry a current during the normal operation of a device. Functional earth connections may be required by devices such as surge suppression and electromagnetic interference filters, some types of antennas and various measurement instruments. Generally the protective earth is also used as a functional earth, though this requires care in some situations.
In household wiring:-
There are two main approaches to the problem of how to disconnect power when a live wire comes into contact with metalwork attached to the earthing system: One way is to get the resistance through the fault path and back to the supply very low by having a metallic connection from the earth back to the supply transformer (a TN system). Then when a fault happens a very high current will flow rapidly blowing a fuse (or tripping a circuit breaker).
The second approach (), where such a direct connection is not used (a TT system), the resistance of the fault path back to the supply is too high for the branch circuit overcurrent protection to operate (blow a fuse or trip a circuit breaker). In such case a residual current detector is installed to detect the current leaking to ground and interrupt the circuit.
Cost:-
TN networks save the cost of a low-impedance earth connection at the site of each consumer. Such a connection (a buried metal structure) is required to provide protective earth in IT and TT systems. TN-C networks save the cost of an additional conductor needed for separate N and PE connections. However, to mitigate the risk of broken neutrals, special cable types and lots of connections to earth are needed. TT networks require proper RCD protection.
Safety:-
In TN, an insulation fault is very likely to lead to a high short-circuit current that will trigger an overcurrent circuit-breaker or fuse and disconnect the L conductors. With TT systems, the earth fault loop impedance can be too high to do this, or too high to do it quickly, so an RCD (or formerly ELCB) is usually employed. The provision of a Residual-current device (RCD) or ELCB to ensure safe disconnection makes these installations EEBAD (Earthed Equipotential Bonding and Automatic Disconnection). Many 1950s and earlier earlier TT installations in the UK may lack this important safety feature. Non-EEBAD installations are capable of the whole installation CPC (Circuit Protective Conductor) remaining live for extended periods under fault conditions, which is a real danger. In TN-S and TT systems (and in TN-C-S beyond the point of the split), a residual-current device can be used as an additional protection. In the absence of any insulation fault in the consumer device, the equation IL1+IL2+IL3+IN = 0 holds, and an RCD can disconnect the supply as soon as this sum reaches a threshold (typically 10-500 mA). An insulation fault between either L or N and PE will trigger an RCD with high probability. In IT and TN-C networks, residual current devices are far less likely to detect an insulation fault. In a TN-C system, they would also be very vulnerable to unwanted triggering from contact between earth conductors of circuits on different RCDs or with real ground, thus making their use impracticable. Also, RCDs usually isolate the neutral core. Since it is unsafe to do this in a TN-C system, RCDs on TN-C should be wired to only interrupt the live conductor. In single-ended single-phase systems where the Earth and neutral are combined (TN-C, and the part of TN-C-S systems which uses a combined neutral and earth core), if there is a contact problem in the PEN conductor, then all parts of the earthing system beyond the break will rise to the potential of the L conductor. In an unbalanced multi-phase system, the potential of the earthing system will move towards that of the most loaded live conductor. Therefore, TN-C connections must not go across plug/socket connections or flexible cables, where there is a higher probability of contact problems than with fixed wiring. There is also a risk if a cable is damaged, which can be mitigated by the use of concentric cable construction and/or multiple earth electrodes. Due to the (small) risks of the lost neutral, use of TN-C-S supplies is banned for caravans and boats in the UK, and it is often recommended to make outdoor wiring TT with a separate earth electrode. In IT systems, a single insulation fault is unlikely to cause dangerous currents to flow through a human body in contact with earth, because no low-impedance circuit exists for such a current to flow. However, a first insulation fault can effectively turn an IT system into a TN system, and then a second insulation fault can lead to dangerous body currents. Worse, in a multi-phase system, if one of the live conductors made contact with earth, it would cause the other phase cores to rise to the phase-phase voltage relative to earth rather than the phase-neutral voltage. IT systems also experience larger transient overvoltages than other systems. In TN-C and TN-C-S systems, any connection between the combined neutral-and-earth core and the body of the earth could end up carrying significant current under normal conditions, and could carry even more under a broken neutral situation. Therefore, main equipotential bonding conductors must be sized with this in mind; use of TN-C-S is inadvisable in situations such as petrol stations, where there is a combination of lots of buried metalwork and explosive gases.
Electromagnetic compatibility:-
In TN-S and TT systems, the consumer has a low-noise connection to earth, which does not suffer from the voltage that appears on the N conductor as a result of the return currents and the impedance of that conductor. This is of particular importance with some types of telecommunication and measurement equipment. In TT systems, each consumer has its own high-quality connection with earth, and will not notice any currents that may be caused by other consumers on a shared PE line.
Regulations:-

In residential and commercial installations in the U. S. and Canada, the feed from the distribution transformer uses a combined neutral and grounding conductor (two phase and one neutral, for three wires total), but within the structure separate neutral and protective earth conductors are used (TN-C-S). The neutral must be connected to the earth (ground) conductor only on the supply side of the customer's disconnecting switch. Additional connections of neutral to ground within the customer's wiring are prohibited. For wiring less than 1000 V, the United States National Electrical Code and Canadian Electrical Code forbid the use of systems that combine the grounding conductor and neutral beyond the customer's disconnecting switch. Exemptions are made for certain appliances, such as: electric clothes dryers, and electric stoves. In Argentina, France (TT) and Australia (TN-C-S), the customer must provide their own ground connection. Japan is governed by PSE law.
Application examples:-

Most modern homes in Europe have a TN-C-S earthing system. The combined neutral and earth occurs between the nearest transformer substation and the service cut out (the fuse before the meter). After this, separate earth and neutral cores are used in all the internal wiring. Older urban and suburban homes in the UK tend to have TN-S supplies, with the earth connection delivered through the lead sheath of the underground lead-and-paper cable. Some older homes, especially those built before the invention of residual-current circuit breakers and wired home area networks, use an in-house TN-C arrangement. This is no longer recommended practice. Laboratory rooms, medical facilities, construction sites, repair workshops, mobile electrical installations, and other environments that are supplied via engine-generators where there is an increased risk of insulation faults, often use an IT earthing arrangement supplied from isolation transformers. To mitigate the two-fault issues with IT systems, the isolation transformers should supply only a small number of loads each and/or should be protected with an insulation monitoring device (generally used only by medical, railway or military IT systems, because of cost). In remote areas, where the cost of an additional PE conductor outweighs the cost of a local earth connection, TT networks are commonly used in some countries, especially in older properties. TT supplies to individual properties are also seen in mostly TN-C-S systems where an individual property is considered unsuitable for TN-C-S supply. In Australia, the TN-C-S system is in use; however, the wiring rules currently state that, in addition, each customer must provide a separate connection to earth via both a water pipe bond (if metallic water pipes enter the consumer's premises) and a dedicated earth electrode. In older installations, it is not uncommon to find only the water pipe bond, and it is allowed to remain as such, but the additional earth electrode must be installed if any upgrade work is done. The protective earth and neutral conductors are combined until the consumer's neutral link (located on the customer's side of the electricity meter's neutral connection) - beyond this point, the protective earth and neutral conductors are separate.

voltage regulator

2009-08-22T10:22:00.000-07:00

A voltage regulator is an electrical regulator designed to automatically maintain a constant voltage level.
It may use an electromechanical mechanism, or passive or active electronic components. Depending on the design, it may be used to regulate one or more AC or DC voltages.
With the exception of passive shunt regulators, all modern electronic voltage regulators operate by comparing the actual output voltage to some internal fixed reference voltage. Any difference is amplified and used to control the regulation element in such a way as to reduce the voltage error. This forms a negative feedback servo control loop; increasing the open-loop gain tends to increase regulation accuracy but reduce stability (avoidance of oscillation, or ringing during step changes). There will also be a trade-off between stability and the speed of the response to changes. If the output voltage is too low (perhaps due to input voltage reducing or load current increasing), the regulation element is commanded, up to a point, to produce a higher output voltage - by dropping less of the input voltage (for linear series regulators and buck switching regulators), or to draw input current for longer periods (boost-type switching regulators); if the output voltage is too high, the regulation element will normally be commanded to produce a lower voltage. However, many regulators have over-current protection, so entirely stop sourcing current (or limit the current in some way) if the output current is too high, and some regulators may also shut down if the input voltage is outside a given range (see also: crowbar circuits).
Measures of regulator quality:-
The output voltage can only be held roughly constant; the regulation is specified by two measurements:
load regulation is the change in output voltage for a given change in load current (for example: "typically 15mV, maximum 100mV for load currents between 5mA and 1.4A, at some specified temperature and input voltage"). line regulation or input regulation is the degree to which output voltage changes with input (supply) voltage changes - as a ratio of output to input change (for example "typically 13mV/V"), or the output voltage change over the entire specified input voltage range (for example "plus or minus 2% for input voltages between 90V and 260V, 50-60Hz"). Other important parameters are:
Temperature coefficient of the output voltage is the change in output voltage with temperature (perhaps averaged over a given temperature range), while... Initial accuracy of a voltage regulator (or simply "the voltage accuracy") reflects the error in output voltage for a fixed regulator without taking into account temperature or aging effects on output accuracy. Dropout voltage - the minimum difference between input voltage and output voltage for which the regulator can still supply the specified current. A Low Drop-Out (LDO) regulator is designed to work well even with an input supply only a Volt or so above the output voltage. Absolute Maximum Ratings are defined for regulator components, specifying the continuous and peak output currents that may be used (sometimes internally limited), the maximum input voltage, maximum power dissipation at a given temperature, etc. Output noise (thermal white noise) and output dynamic impedance may be specified as graphs versus frequency, while output ripple noise (mains "hum" or switch-mode "hash" noise) may be given as peak-to-peak or RMS voltages, or in terms of their spectra. Quiescent current in a regulator circuit is the current drawn internally, not available to the load, normally measured as the input current while no load is connected (and hence a source of inefficiency; some linear regulators are, surprisingly, more efficient at very low current loads than switch-mode designs because of this).
Electromechanical regulators:-
In older electromechanical regulators, voltage regulation is easily accomplished by coiling the sensing wire to make an electromagnet. The magnetic field produced by the current attracts a moving ferrous core held back under spring tension or gravitational pull. As voltage increases, so does the current, strengthening the magnetic field produced by the coil and pulling the core towards the field. The magnet is physically connected to a mechanical power switch, which opens as the magnet moves into the field. As voltage decreases, so does the current, releasing spring tension or the weight of the core and causing it to retract. This closes the switch and allows the power to flow once more.
If the mechanical regulator design is sensitive to small voltage fluctuations, the motion of the solenoid core can be used to move a selector switch across a range of resistances or transformer windings to gradually step the output voltage up or down, or to rotate the position of a moving-coil AC regulator.
Early automobile generators and alternators had a mechanical voltage regulator using one, two, or three relays and various resistors to stabilize the generator's output at slightly more than 6 or 12 V, independent of the engine's rpm or the varying load on the vehicle's electrical system. Essentially, the relay(s) employed pulse width modulation to regulate the output of the generator, controlling the field current reaching the generator (or alternator) and in this way controlling the output voltage produced.
The regulators used for generators (but not alternators) also disconnect the generator when it was not producing electricity, thereby preventing the battery from discharging back through the stopped generator. The rectifier diodes in an alternator automatically perform this function so that a specific relay is not required; this appreciably simplified the regulator design.
More modern designs now use solid state technology (transistors) to perform the same function that the relays perform in electromechanical regulators.
Main regulator:-
Electromechanical regulators have also been used to regulate the voltage on AC power distribution lines. These regulators generally operate by selecting the appropriate tap on a transformer with multiple taps. If the output voltage is too low, the tap changer switches connections to produce a higher voltage. If the output voltage is too high, the tap changer switches connections to produce a lower voltage. The controls provide a deadband wherein the controller will not act, preventing the controller from constantly hunting (constantly adjusting the voltage) to reach the desired target voltage.
Coil-rotation AC voltage regulator:-
This is an older type of regulator used in the 1920s that uses the principle of a fixed-position field coil and a second field coil that can be rotated on an axis in parallel with the fixed coil.
When the movable coil is positioned perpendicular to the fixed coil, the magnetic forces acting on the movable coil balance each other out and voltage output is unchanged. Rotating the coil in one direction or the other away from the center position will increase or decrease voltage in the secondary movable coil.
This type of regulator can be automated via a servo control mechanism to advance the movable coil position in order to provide voltage increase or decrease. A braking mechanism or high ratio gearing is used to hold the rotating coil in place against the powerful magnetic forces acting on the moving coil.
AC voltage stabilizer:-
A voltage stabilizer is a type of household mains regulator which uses a continuously variable autotransformer to maintain an AC output that is as close to the standard or normal mains voltage as possible, under conditions of fluctuation. It uses a servomechanism (or negative feedback) to control the position of the tap (or wiper) of the autotransformer, usually with a motor. An increase in the mains voltage causes the output to increase, which in turn causes the tap (or wiper) to move in the direction that reduces the output towards the nominal voltage.
An alternative method is the use of a type of saturating transformer called a ferroresonant transformer or constant-voltage transformer. These transformers use a tank circuit composed of a high-voltage resonant winding and a capacitor to produce a nearly constant average output with a varying input. The ferroresonant approach is attractive due to its lack of active components, relying on the square loop saturation characteristics of the tank circuit to absorb variations in average input voltage. Older designs of ferroresonant transformers had an output with high harmonic content, leading to a distorted output waveform. Modern devices are used to construct a perfect sinewave. The ferroresonant action is a flux limiter rather than a voltage regulator, but with a fixed supply frequency it can maintain an almost constant average output voltage even as the input voltage varies widely.
The ferroresonant transformers, which are also known as Constant Voltage Transformers (CVTs) or ferros, are also good surge suppressors, as they provide high isolation and inherent shortcircuit protections.
A ferroresonant transformer can operate with an input voltage range ±40% or more of the nominal voltage.
Output power factor remains in the range of 0.96 or higher from half to full load.
Because it regenerates an output voltage waveform, output distortion, which is typically less than 4%, is independent of any input voltage distortion, including notching.
Efficiency at full load is typically in the range of 89% to 93%. However, at low loads, efficiency can drop below 60% and no load losses can be as high as 20%. The current-limiting capability also becomes a handicap when a CVT is used in an application with moderate to high inrush current like motors, transformers or magnets. In this case, the CVT has to be sized to accommodate the peak current, thus forcing it to run at low loads and poor efficiency.
Minimum maintenance is required beyond annual replacement of failed capacitors. Redundant capacitors built into the units allow several capacitors to fail between inspections without any noticeable effect to the device's performance.
Output voltage varies about 1.2% for every 1% change in supply frequency. For example, a 2-Hz change in generator frequency, which is very large, results in an output voltage change of only 4%, which has little effect for most loads.
It accepts 100% single-phase switch-mode power supply loading without any requirement for derating, including all neutral components.
Input current distortion remains less than 8% THD even when supplying nonlinear loads with more than 100% current THD.
Drawbacks of CVTs (constant voltage transformers) are their larger size, high audible humming sound, and high heat generation
DC voltagr stsbilizers:-

Many simple DC power supplies regulate the voltage using a shunt regulator such as a zener diode, avalanche breakdown diode, or voltage regulator tube. Each of these devices begins conducting at a specified voltage and will conduct as much current as required to hold its terminal voltage to that specified voltage. The power supply is designed to only supply a maximum amount of current that is within the safe operating capability of the shunt regulating device (commonly, by using a series resistor). In shunt regulators, the voltage reference is also the regulating device.
If the stabilizer must provide more power, the shunt regulator output is only used to provide the standard voltage reference for the electronic device, known as the voltage stabilizer. The voltage stabilizer is the electronic device, able to deliver much larger currents on demand.

Sprinkler

2009-08-22T10:16:00.000-07:00

Sprinkler:-
Irrigation sprinklers, a device for irrigation of lawns or crops Fire sprinkler system, the entire systems of pipes and sprinklers intended for fire suppression within buildings Fire sprinkler, a device for fire suppression within buildings Sprinkler (band), a rock band.

Smoke alarm

2009-08-21T02:59:00.000-07:00

A smoke detector is a device that detects smoke. Commercial, industrial, and mass residential devices issue a signal to a fire alarm system, while household detectors, known as smoke alarms, generally issue a local audible and/or visual alarm from the detector itself.
Smoke detectors are typically housed in a disk-shaped plastic enclosure about 150 millimetres (6 in) in diameter and 25 millimetres (1 in) thick, but the shape can vary by manufacturer or product line. Most smoke detectors work either by optical detection (photoelectric) or by physical process (ionization), while others use both detection methods to increase sensitivity to smoke. Smoke detectors in large commercial, industrial, and residential buildings are usually powered by a central fire alarm system, which is powered by the building power with a battery backup. However, in many single family detached and smaller multiple family housings, a smoke alarm is often powered only by a single disposable battery.
History:-
The first automatic electric fire alarm was invented in 1890 by Francis Robbins Upton (US patent no. 436,961). Upton was an associate of Thomas Edison, although there is no evidence that Edison contributed to this project.
In the late 1930s the Swiss physicist Walter Jaeger tried to invent a sensor for poison gas. He expected that gas entering the sensor would bind to ionized air molecules and thereby alter an electric current in a circuit in the instrument. His device failed: small concentrations of gas had no effect on the sensor's conductivity. Frustrated, Jaeger lit a cigarette—and was soon surprised to notice that a meter on the instrument had registered a drop in current. Smoke particles had apparently done what poison gas could not. Jaeger's experiment was one of the advances that paved the way for the modern smoke detector.
It was 30 years, however, before progress in nuclear chemistry and solid-state electronics made a cheap sensor possible. While home smoke detectors were available during most of the 1960s, the price of these devices was rather high. Before that, alarms were so expensive that only major businesses and theaters could afford them.
The first truly affordable home smoke detectors were invented by Duane D. Pearsall and Stanley Bennett Peterson in 1965, featuring individual battery powered units that could be easily installed and replaced. The first units for mass production came from the manufacturing mind of Stanley B. Peterson in 1975 at Duane Pearsall’s company in Lakewood, Colorado, named Statitrol Corporation.
These first units were made from strong fire resistant steel and shaped much like a bee's hive. The battery was a rechargeable specialized unit created by Gates Energy. The need of a quick replace battery didn't take long to show its need and the rechargeable was replaced with a pair of double A's along with a plastic shell encasing the detector. Daniel Peterson, son to Stanley, was the manufacturing floor supervisor and shipping lead for quality control setting a new record for small business PCB electro/mechanical manufacturing during the time (1975–1977). The small assembly line sent close to 500 units per day before Statitrol sold its invention to Emerson Electric in 1980 and Sears’s retailers picked up full distribution of the 'now required in every home' smoke detector.
The first commercial smoke detectors came to market in 1969. Today they are installed in 93% of US homes and 85% of UK homes. However it is estimated that any given time over 30% of these alarms don't work, as users remove the batteries, or forget to replace them.
Although commonly attributed to NASA, smoke detectors were not invented as a result of the space program, though a variant with adjustable sensitivity was developed for Skylab
Desine:-
An optical detector is a light sensor. When used as a smoke detector, it includes a light source (incandescent bulb or infrared LED), a lens to collimate the light into a beam, and a photodiode or other photoelectric sensor at an angle to the beam as a light detector. In the absence of smoke, the light passes in front of the detector in a straight line. When smoke enters the optical chamber across the path of the light beam, some light is scattered by the smoke particles, directing it at the sensor and thus triggering the alarm.
Also seen in large rooms, such as a gymnasium or an auditorium, are devices to detect a projected beam. A unit on the wall sends out a beam, which is either received by a receiver or reflected back via a mirror. When the beam is less visible to the "eye" of the sensor, it sends an alarm signal to the fire alarm control panel.
Optical smoke detectors are quick in detecting particulate (smoke) generated by smoldering (cool, smoky) fires. Many independent tests indicate that optical smoke detectors typically detect particulates (smoke) from hot, flaming fires approximately 30 seconds later than ionization smoke alarms.
They are less sensitive to false alarms from steam or cooking fumes generated in kitchen or steam from the bathroom than are ionization smoke alarms. For the aforementioned reason, they are often referred to as 'toast proof' smoke alarms.
Lonization:-
This type of detector is cheaper than the optical detector; however, it is sometimes rejected because it is more prone to false alarms than photoelectric smoke detectors[citation needed]. It can detect particles of smoke that are too small to be visible. It includes 150 milligrams of radioactive americium 241 (241Am). The radiation passes through an ionization chamber, an air-filled space between two electrodes, and permits a small, constant current between the electrodes. Any smoke that enters the chamber absorbs the alpha particles, which reduces the ionization and interrupts this current, setting off the alarm.
241Am, an alpha emitter, has a half-life of 432.2 years. This means that it does not have to be replaced during the useful life of the detector, and also makes it safe for people at home, since it is only slightly radioactive. Alpha radiation, as opposed to beta and gamma, is used for two additional reasons: Alpha particles have high ionization, so sufficient air particles will be ionized for the current to exist, and they have low penetrative power, meaning they will be stopped by the plastic of the smoke detector and/or the air.
Air-sampling:-
An air-sampling smoke detector is capable of detecting microscopic particles of smoke. Most air-sampling detectors are aspirating smoke detectors, which work by actively drawing air through a network of small-bore pipes laid out above or below a ceiling in parallel runs covering a protected area. Small holes drilled into each pipe form a matrix of holes (sampling points), providing an even distribution across the pipe network. Air samples are drawn past a sensitive optical device, often a solid-state laser, tuned to detect the extremely small particles of combustion. Air-sampling detectors may be used to trigger an automatic fire response, such as a gaseous fire suppression system, in high-value or mission-critical areas, such as archives or computer server rooms.
Most air-sampling smoke detection systems are capable of a higher sensitivity than spot type smoke detectors and provide multiple levels of alarm threshold, such as Alert, Action, Fire 1 and Fire 2. Thresholds may be set at levels across a wide range of smoke levels. This provides earlier notification of a developing fire than spot type smoke detection, allowing manual intervention or activation of automatic suppression systems before a fire has developed beyond the smoldering stage, thereby increasing the time available for evacuation and minimizing fire damage.
Carbon monoxide and carbon dioxide detection:-
Some smoke alarms use a carbon dioxide sensor or carbon monoxide sensor in order to detect extremely dangerous products of combustion. However, not all smoke detectors that are advertised with such gas sensors are actually able to warn of poisonous levels of those gases in the absence of a fire.
Performance differences:-
Optical or "toast-proof" smoke detectors are generally quicker in detecting particulate (smoke) generated by smoldering (cool, smokey) fires. Ionization smoke detectors are generally quicker in detecting particulate (smoke) generated by flaming (hot) fires.
According to fire tests conformant to EN 54, normally the CO2 cloud from smoke can be detected before particulate.
Obscuration is a unit of measurement that has become the standard definition of smoke detector sensitivity. Obscuration is the effect that smoke has on reducing visibility. Higher concentrations of smoke result in higher obscuration levels, lowering visibility.

Fire alarm system

2009-08-21T02:31:00.000-07:00

An automatic fire alarm system is designed to detect the unwanted presence of fire by monitoring environmental changes associated with combustion. In general, a fire alarm system is either classified as automatic, manually activated, or both. Automatic fire alarm systems can be used to notify people to evacuate in the event of a fire or other emergency, to summon emergency services, and to prepare the structure and associated systems to control the spread of fire and smoke.
Design:-

After the fire protection goals are established - usually by referencing the minimum levels of protection mandated by the appropriate model building code, insurance agencies, and other authorities - the fire alarm designer undertakes to detail specific components, arrangements, and interfaces necessary to accomplish these goals. Equipment specifically manufactured for these purposes are selected and standardized installation methods are anticipated during the design. In the United States, NFPA 72, The National Fire Alarm Code is an established and widely used installation standard.
Fundamental configuration:-
Fire alarm control panel: This component, the hub of the system, monitors inputs and system integrity, control outputs and relays information. Primary Power supply: Commonly the non-switched 120 or 240 V AC source supplied from a commercial power utility. In non-residential applications, a branch circuit is dedicated to the fire alarm system and its constituents. "Dedicated branch circuits" should not be confused with "Individual branch circuits" which supply energy to a single appliance. Secondary (backup) Power supplies: This component, commonly consisting of sealed lead-acid storage batteries or other emergency sources including generators, is used to supply energy in the event of a primary power failure. Initiating Devices: This component acts as input to the fire alarm control unit and are either manually or automatically activated. Notification appliances: This component uses energy supplied from the fire alarm system or other stored energy source, to inform the proximate persons of the need to take action, usually to evacuate. Building Safety Interfaces: This interface allows the fire alarm system to control aspects of the built environment and to prepare the building for fire and to control the spread of smoke fumes and fire by influencing air movement, lighting, process control, human transport and exit.
Initiating devices:-

Manually activated devices; Break glass stations, Buttons and manual pull station are constructed to be readily located (near the exits), identified, and operated. Automatically activated devices can take many forms intended to respond to any number of detectable physical changes associated with fire: convected thermal energy; heat detector, products of combustion; smoke detector, radiant energy; flame detector, combustion gasses; carbon monoxide detector and release of extinguishing agents; water-flow detector. The newest innovations can use cameras and computer algorithms to analyze the visible effects of fire and movement in applications inappropriate for or hostile to other detection methods.
Notification appliances:-
Audible, visible, tactile, textual or even olfactory stimuli (odorizer). to alert the occupants. Audible or visible signals are the most common and may utilize speakers to deliver live or pre-recorded instructions to the occupants. In the United States, fire alarm evacuation signals are required to use a standardized interrupted four count temporal pattern to avoid confusion with other signals using similar sounding appliances. Other methods include: Audible textual appliances, which are employed as part of a fire alarm system that includes Emergency Voice Alarm Communications (EVAC) capabilities. High reliability speakers are used to notify the occupants of the need for action in connection with a fire or other emergency. These speakers are employed in large facilities where general undirected evacuation is considered impracticable or undesirable. The signals from the speakers are used to direct the occupant's response. The system may be controlled from one or more locations within the building known as Fire Wardens Stations, or from a single location designated as the building Fire Command Center. Speakers are automatically activated by the fire alarm system in a fire event, and following a pre-alert tone, selected groups of speakers may transmit one or more prerecorded messages directing the occupants to safety. These messages may be repeated in one or more languages. Trained personnel activating and speaking into a dedicated microphone can suppress the replay of automated messages in order to initiate or relay real time voice instructions.
Building safety interfaces:-
Magnetic Smoke Door Holders: Wall or floor mounted solenoids or electromagnets controlled by a fire alarm system or detection component that magnetically secures spring-loaded self-closing smoke tight doors in the open position. Designed to de-magnetize to allow automatic closure of the door on command from the fire control or upon failure of the power source, interconnection or controlling element. Stored energy in the form of a spring or gravity can then close the door to restrict the passage of smoke from one space to another in an effort to maintain a tenable atmosphere on either side of the door during evacuation and fire fighting efforts. Duct Mounted Smoke Detection: Smoke detection mounted in such a manner as to sample the airflow through duct work and other plenums specifically fabricated for the transport of environmental air into conditioned spaces. Interconnection to the fan motor control circuits are intended to stop air movement, close dampers and generally prevent the recirculation of toxic smoke and fumes produced by fire into occupiable spaces.
UK fire alarm system categories:-

There are many types of fire alarm systems each suited to different building types and applications. A fire alarm system can vary dramatically in both price and complexity, from a single panel with a detector and sounder in a small commercial property to an addressable fire alarm system in a multi-occupancy building. Systems have to protect both buildings and occupants.
The categories of fire alarm systems are L if they are designed to protect life, P to protect buildings and M if they are manual systems.
M Manual systems, e.g. hand bells, gongs, etc. These may be purely manual or manual electric, the latter may have call points and sounders. They rely on the occupants of the building discovering the fire and acting to warn others by operating the system. Such systems form the basic requirement for places of employment with no sleeping risk. P1 The system is installed throughout the building - the objective being to call the fire brigade as early as possible to ensure that any damage caused by fire is minimized. Small low risk areas can be excepted, such as toilets and cupboards less than 1m². P2 Detection should be provided in parts of the building where the risk of ignition is high and/or the contents are particularly valuable. Category 2 systems provide fire detection in specified parts of the building where there is either high risk or where business disruption must be minimised. L1 A category L1 system is designed for the protection of life and which has automatic detectors installed throughout all areas of the building (including roof spaces and voids) with the aim of providing the earliest possible warning. A category L1 system is likely to be appropriate for the majority of residential care premises. In practice, detectors should be placed in nearly all spaces and voids. With category 1 systems, the whole of a building is covered apart from minor exceptions. L2 A category L2 system designed for the protection of life and which has automatic detectors installed in escape routes, rooms adjoining escape routes and high hazard rooms. In a medium sized premises (sleeping no more than ten residents), a category L2 system is ideal. These fire alarm systems are identical to an L3 system but with additional detection in an area where there is a high chance of ignition, e.g., kitchen) or where the risk to people is particularly increased (e.g., sleeping risk). L3 This category is designed to give early warning to everyone. Detectors should be placed in all escape routes and all rooms that open onto escape routes. Category 3 systems provide more extensive cover than category 4. The objective is to warn the occupants of the building early enough to ensure that all are able to exit the building before escape routes become impassable. L4 Category 4 systems cover escape routes and circulation areas only. Therefore, detectors will be placed in escape routes, although this may not be suitable depending on the risk assessment or if the size and complexity of a building is increased. Detectors might be sited in other areas of the building, but the objective is to protect the escape route. L5 This is the "all other situations" category, e.g., computer rooms, which may be protected with an extinguishing system triggered by automatic detection. Category 5 systems are the "custom" category and relate to some special requirement that cannot be covered by any other category.

Float switch

2009-08-20T08:28:00.000-07:00

A float switch is a device used to detect the level of liquid within a tank. The switch may be used in a pump, an indicator, an alarm, or other devices.
Float switches range from small to large and may be as simple as a mercury switch inside a hinged float or as complex as a series of optical or conductance sensors producing discrete outputs as the liquid reaches many different levels within the tank. Perhaps the most common type of float switch is simply a float raising a rod that actuates a microswitch.
A very common application is in sump pumps and condensate pumps where the switch detects the rising level of liquid in the sump or tank and energizes an electrical pump which then pumps liquid out until the level of the liquid has been substantially reduced, at which point the pump is switched off again. Float switches are often adjustable and can include substantial hysteresis. That is, the switch's "turn on" point may be much higher than the "shut off" point. This minimizes the on-off cycling of the associated pump.
Some float switches contain a two-stage switch. As liquid rises to the trigger point of the first stage, the associated pump is activated. If the liquid continues to rise (perhaps because the pump has failed or its discharge is blocked), the second stage will be triggered. This stage may switch off the source of the liquid being pumped, trigger an alarm, or both.